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Techniques and Applications of Path Integration
Techniques and Applications of Path Integration
Techniques and Applications of Path Integration
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Techniques and Applications of Path Integration

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A book of techniques and applications, this text defines the path integral and illustrates its uses by example. It is suitable for advanced undergraduates and graduate students in physics; its sole prerequisite is a first course in quantum mechanics. For applications requiring specialized knowledge, the author supplies background material.
The first part of the book develops the techniques of path integration. Topics include probability amplitudes for paths and the correspondence limit for the path integral; vector potentials; the Ito integral and gauge transformations; free particle and quadratic Lagrangians; properties of Green's functions and the Feynman-Kac formula; functional derivatives and commutation relations; Brownian motion and the Wiener integral; and perturbation theory and Feynman diagrams.
The second part, dealing with applications, covers asymptotic analysis and the calculus of variations; the WKB approximation and near caustics; the phase of the semiclassical amplitude; scattering theory; and geometrical optics. Additional topics include the polaron; path integrals for multiply connected spaces; quantum mechanics on curved spaces; relativistic propagators and black holes; applications to statistical mechanics; systems with random impurities; instantons and metastability; renormalization and scaling for critical phenomena; and the phase space path integral.
LanguageEnglish
Release dateOct 10, 2012
ISBN9780486137025
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    Techniques and Applications of Path Integration - L. S. Schulman

    Supplements

    PART ONE

    Introduction

    ONE

    Introducing and Defining the Path Integral

    The best place to find out about path integrals is in Feynman’s paper.a Our approach is not to use path integrals as a way of arriving at quantum mechanics, although Feynman has used this point of view in his book with Hibbs. Rather we assume knowledge of quantum mechanics and deduce the path integral formalism from it. This gets us into the subject quickly.

    The wave function of a nonrelativistic spinless particle in one dimension evolves according to Schrödinger’s equation

    (1.1)

    (1.2)

    Our interest is in the propagator or Green’s function G which satisfies the equation

    (1.3)

    in operator notation. In coordinate space this is written

    (1.4)

    The G’s are related by

    (1.5)

    Knowing G means having a solution to the time dependent Schrödinger equation in the sense that if ψ(t0) is the state of the system at t0, ψ(t), given by

    (1.6)

    is the state at t. For time independent H an operator solution of (1.3) can immediately be written down:

    (1.7)

    where θ is the step function. Since H is assumed to be time independent we can, without loss of generality, take t0=0. Then for t>0 we have

    (1.8)

    where the argument 0 has been deleted.

    The path integral arises from the fact that

    (1.9)

    Letting λ = it/ yields

    (1.10)

    with the product in the brackets taken N times. Now we make use of a fundamental fact about the exponential of two operators, namely

    (1.11)

    This is proved easily enough,b and in a power series expansion the coefficient of the λ²/N² term is

    In subsequent manipulations we assume that the O(1/N²) term is well behaved, that it stays bounded when applied to states, and so on. For reasonable potentials this assumption is justified; more is said on this topic in the appendix.

    What we are now aiming for is to replace the term

    (1.12)

    by the term

    (1.13)

    For real numbers (rather than operators) this replacement is a reflection of a fundamental fact about the exponential. The expression

    converges to ex despite the presence of yn so long as yn→0 as n→∞. (A proof of this assertion can be had by taking large enough n that |yn| <δ and by using the bound

    and assuming n large enough that |(x-δ)/n| < 1.)

    For operators a bit of care is required, and the trick is to express the difference of (1.12) and (1.13) in a peculiar way:

    (1.14)

    Equation 1.14 is an identity. It containsN terms, each of which has the factor exp(-λT/N)exp(-λV/N)-exp(-λ(T+V)/N), which by (1.11) is of order 1/N². Hence in the limit the difference is zero. (In an appendix mention is made of various finer points in the estimate.)

    We have therefore justified the replacement of (1.10) by

    (1.15)

    In effect we have given a heuristic proof of the Trotter product formula. From here, getting the path integral is just a few easy steps. The identity operator, in the form

    (1.16)

    is inserted between each term in the product in (1.15), yielding

    (1.17)

    (for convenience we have taken y = x0, x = xN,). The multiplication operator V is diagonal in coordinate space so that

    (1.18)

    Next we require coordinate space matrix elements of exp(-λT/N) (between states 〈η| and |ξ〉, say), and to obtain these we insert a complete set of momentum states

    (1.19)

    This gives

    (1.20)

    This is our first Gaussian integral of the book, but far from the last. The general formula is

    (1.21)

    Using (1.21), (1.20) becomes

    (1.22)

    Equations 1.18 and 1.22 are inserted in (1.17) to yield

    (1.23)

    Now let ε = t/Nλ/iN and combine the exponentials in (1.23):

    (1.24)

    Equation 1.24 is the path integral expression for the propagator. A few words are in order, however, on why this is called a path integral or sum over histories.

    Imagine that the points y, x1,... xN-1, x are connected by lines. Then we have a broken line path from y to x. The sum in the exponential of (1.24) can be interpreted as a Riemann sum of a certain integral along that path:

    (1.25)

    The integrand in (1.25) is well known in classical mechanics. It is just the Lagrangian

    (1.26)

    of the classical system, which when quantized has the Hamiltonian (1.2). Furthermore, the action

    (1.27)

    is no less prominent an object in classical mechanics. The argument of the exponential in (1.24) is thus iS, with S evaluated along the broken line path connecting y, x1, x2,... xN-1, x.

    The integrals over the quantities x1,...xN-1 can be interpreted as summing over all possible broken line paths connecting y and x. Since any continuous path can be approximated by a broken line path, and considering the fact that the limit N→∞ is taken, one might be optimistic enough to interpret the integrals as a sum over all paths. A final cosmetic transformation on (1.24) is to let

    (1.28)

    and call this a normalization constant which, though infinite in the limit N→∞, serves merely to make sure G is a unitary operator (more about this later). Equation 1.24 has become

    (1.29)

    This is a sum over paths, or histories, of eiS[xwith all paths satisfying x(0)=y, x(t)=x, entering the sum. The capital Σ is used to avoid giving the impression that we have a bona fide measure, and the present state of knowledge on path integration is such that when you encounter expressions like (1.29) they mean neither more nor less than what appears in (1.24)—despite the more suggestive notation some authors prefer to use.

    Exercise: Derive a path integral expression for the Green’s function in the case of a time dependent Hamiltonian of the form

    HINTS:

    (1) Although (1.7) is no longer true, we still have

    (1.30)

    with (t1 - tt factors. Equation 1.30 is just a time ordered product, about which more later.

    (2) It may be useful to change time variables.

    APPENDIX: THE TROTTER PRODUCT FORMULA

    THEOREM: (Trotter product formula) Let A and B be linear operators on a Banach space X such that A, B and A+B are the infinitesimal generators of the contraction semigroups Pt, Qt, and Rt respectively. Then for all ψ∈X

    In this appendix we define some of the terms used above, indicate how a proof of the theorem goes, and examine some of its consequences.

    A semigroup is a set closed under a binary, associative operation. Were inverses required to be in the set it would form a group. A semigroup may or may not possess an identity.

    Definition: A contraction semigroup on Banach space X is a family of bounded everywhere defined linear operators Pt, 0 ≤ t < ∞ mapping X→X such that

    The norm used above is defined as follows:

    and ||x|| is the norm in X. The term contraction comes from the fact that ||Pt||≤1, since vectors do not grow as they evolve under Pt. The infinitesimal generator A of Pt is defined by

    on the domain D(A) of all ψ∈X for which the limit exists.

    Remarks on the Proof of the Theorem (leaving out most statements about domains): let h be a positive real number and let P, Q, R be as defined in the statement of the theorem. By the definition of generators we have

    (PhQh-1)ψ=(Ph-1)ψ+Ph(Qh-1)ψ=h(A+B)ψ+o(h)

    where o(h) denotes vectors x ||x||/h=0. Then since

    (Rh-1)ψ=h(A+B)ψ+o(h)

    it follows that

    (PhQh-Rh)ψ=o(h).

    Now we must establish the uniformity of the bound o(h). By using properties necessarily possessed by infinitesimal generators we show that for ψ in some compact subset of D(A + B), h-1||( PhQh - Rh)ψ|| is uniformly bounded. For some ψ∈D(A+B), {Rsψ}, 0≤st, is compact and in D(A+B), hence ||(PhQh - Rh)Rsψ||=o(h) uniformly in s.

    Let h=t/n. Then we wish to show that

    ||((PhQh)n-Rhn)ψ||→0 for n→∞

    To this end we examine

    Next apply this to ψ, and use the fact that ||PhQh||≤1. This implies

    where the limit is uniform.

    For the physical application of this formula we let

    A = i , B = -iV

    - (A + B) = i+ V) = iH

    Thus it is necessary to know whether A, B, and A + B generate contractive semigroups. Basically what is involved is examining ||eCt|| where C is the proposed infinitesimal generator. Recall that

    For V there is no problem; V is a multiplication operator and

    implies ||e-iVt||=1 for all t.

    To show that the Laplacian Δ generates a contractive semigroup we review some elementary facts and intuitions about the norm. Generally speaking the norm looks for the largest eigenvalue. If X is finite (M) dimensional we have Qψi = qiψi, i = 1,..., M. Then the worst case is the biggest lqi, for then ||Q||= max. Say Q=eC. Its eigenvalues are eCi where Ci are the eigenvalues of C. Then ||eC||=maxi|eCi|=maxieReCi. Thus

    so that the condition for C to be the generator of a contractive semigroup is that ReCi≤0 for all i.

    On Hilbert space, the thing that would be the eigenvector is not always in the space and the definition of the norm as a limit (sup) must be invoked. For example, let M be the multiplication operator by the function exp( - x²). Thus

    Mψ(x) = e-x2ψ(x)

    Clearly ||Mψ(x)|| ≤ ||ψ|| so that ||M|| ≤ 1. Let

    then

    For all finite N

    (See 11.20 for proof of the inequality.) But for large enough N this gets arbitrarily close to 1, hence ||M|| = 1. However, there is no ψ in the Hilbert space such that Mψ=ψ. Thus 1 is not an eigenvalue of M, but it does have some special properties with respect to M: it is in the spectrum. The spectrum is defined as the complement of the resolvent set where the resolvent set of an operator A is the set of λ for which (λ - A)-1 exists.

    Above we had a condition on the eigenvalues of a finite dimensional matrix C so that it generated a contractive semigroup. In a Hilbert space it is most convenient to state the condition in terms of the spectrum. The condition on an operator C is that Re λ ≤0 for λ in the spectrum. If C=iK, in terms of the eigenvalues (or spectrum) ki, of K this means Im ki ≥ 0 for all i. If e-tC is also to be a contractive semigroup—evolution in both directions—we must have

    Im k=0

    That is, K has only real spectrum, a condition guaranteed by the usual requirement that the Hamiltonian be self-adjoint. Thus ||e-it || = 1 which is the statement that free particle propagation is norm preserving. To determine whether A + B + V is self-adjoint (so that its spectrum would be real). This question is what mathematicians call perturbation theory (and physicists never bother to ask). If it is self-adjoint, then the conditions for the Trotter formula are satisfied.

    As shown above the Feynman integral is justified by the foregoing theorem, where by sum over histories is meant the specific way of doing this sum which amounts to the taking of the nth power of (or the n fold iteration of a certain approximation to the Green’s function for) a finite operator. This justification of the path integral loses some of the intuitive appeal of the Feynman’s formulation.

    NOTES

    The original paper on path integrals is

    R. P. Feynman, Rev. Mod. Phys. 20, 367 (1948)

    A general exposition with emphasis on those applications developed by Feynman himself can be found in

    R. P. Feynman and A. R. Hibbs, Quantum Mechanics and Path Integrals, McGraw-Hill, New York, 1965

    The relation of the Trotter product formula to path integration is lucidly presented in

    E. Nelson, J. Math. Phys. 5, 332 (1964)

    and our appendix follows that paper.

    TWO

    Probabilities and Probability Amplitudes for Paths

    When you learn quantum mechanics, you’re told to forget about your naive, classical idea of particles traveling on trajectories. A particle might be here at one time, and there a bit later, but to speak in terms of the path it took from here to there is to invite contradiction and confusion. The chief parable on which to unlearn your classical intuition is the tale of the two slit experiment. An electron passes through either (or both) of two slits on its way to a screen. The screen is a detector for position and an interference pattern can be seen. Any attempt to verify that the electron went through one slit or the other by a localization of the electron will destroy the interference pattern. How can this lesson, well learned, be consistent with the sum over paths that we have just derived?

    The secret lies in the distinction between probability and probability amplitudes. In probability theory there is a rule for conditional probabilities. Let P(a|b) denote the probability of an event a, given that event b occurred, with similar definitions for P(a|c) and P(b|c). Then it is true that

    (2.1)

    where the sum is over all events (or states) b that can occur between c and a. In quantum mechanics there is a different rule. We work with a probability amplitude, say ϕ, which, like P(a|b) depends on two states and satisfies the relation

    (2.2)

    The sum is again over all possible states b. The difference in the two situations is that ϕ is not a probability, and it is only the square of its absolute value that is interpreted as a probability. This profound difference is discussed at length by Feynman.

    Thus although in (2.2) a sum over b is written it would not be correct to say that the system was either in b1 or b2 or since then we ought to be using (2.1) for |ϕab|².

    The same caution applies to the path integral. True, we have found a sum over paths, but there is no assertion that the system followed definite paths with certain probabilities. Rather, we compute probability amplitudes for the paths and sum the amplitudes.

    Exercise: Analyze the two slit experiment in these terms.

    At this point one can turn the entire inquiry around and, starting from (2.2) as a postulate, derive quantum mechanics. Usually in working with (2.2) one is summing over some finite number of spin states, or even some continuum of states. But to get the dynamics of the Schrödinger equation from (2.2) one must try a far more daring sort of sum. Thus we want the amplitude for a (nonrelativistic, spinless) particle to go from y at time 0 to x at time t. We must sum over all intermediate possibilities, that is, our sum must include a contribution for each trajectory from (y,0) to (x, t). For the trajectory x(τ) (x(0) =y, x(t)=x), call the amplitude ϕ(x(τ)). The total amplitude is

    (2.3)

    Exercise: Use this interpretation plus the meaning of the wave function ψ(x) as a probability amplitude for position to show that ϕ(x,t;y,0) is the Green’s function.

    The only remaining question is what to take for ϕ(x(τ)). It appears that it was Dirac who put Feynman on the right track for finding this functional. In an early paper and also in later editions of his quantum mechanics book, Dirac found that the Green’s function (he calls it transformation function) looks something like

    (2.4)

    with S the solution of the Hamilton-Jacobi equation. Now everyone knows that exp(iS; that’s WKB. What Dirac further observed is that exp(iS) is also a good approximation when the time interval over which G is supposed to propagate goes to zero. Consequently, for short times ε, the propagator from y to x is approximated by

    Then, when one sums over intermediate positions for finite time propagators (in the spirit of (2.2)) one gets (1.24).

    This, I believe, was the main precursor to Feynman’s work.

    NOTES

    The early paper by Dirac is

    P. A. M. Dirac, Physikalische Zeitschrift der Sowjetunion, 3, No. 1 (1933)

    This paper, as well as Feynman’s original Rev. Mod. Phys. paper can be found reprinted in

    J. Schwinger, Quantum Electrodynamics, Dover, New York, 1958

    Additional material along these lines can be found in later editions of the famous text

    P. A. M. Dirac, The Principles of Quantum Mechanics, Oxford, London

    In particular see Sections 31 and 32, pp. 121—130 in the fourth (1958) edition.

    In the usual formulation of quantum mechanics ψ(x) is the amplitude to find the particle at x where "x," the position, has an operator corresponding to it and one can in principle localize at x. Given ϕ(x(τ)) of (2.3), for example, one might ask if one could in principle determine (i.e., measure) whether a particle has taken a particular trajectory. A way of doing this is presented in

    Y. Aharonov and M. Vardi, Meaning of an Individual Feynman Path, Phys. Rev. D 21, 2235 (1980)

    THREE

    Correspondence Limit for the Path Integral (Heuristic)

    We have obtained the propagator in the form of a sum

    (3.1)

    where x) is a path starting at y and ending at x. Starting from (3.1) the correspondence limit of quantum mechanics is a wave of the hand away. The way to get the classical limit of (3.1) is the method of stationary phase and the hand waving consists of assuming that what works for a one dimensional integral works for the sum in (3.1).

    In the method of stationary phase one considers integrals of the form

    (3.2)

    What is sought is the dominant contribution to F as λ→∞. Later we have occasion to define this goal more precisely, in terms of order symbols and asymptotic expansions. The incantation used for simplifying (3.2) is as follows. For large λ, the phase of eiλƒ(t) will vary rapidly unless ƒ′=0. This implies that the dominant contribution to the integral comes from regions of t where ƒ′ vanishes. Suppose ƒ′ vanishes at only one point t0. Neglecting contributions to the integral from regions far from t0, ƒ is expanded about t0,

    (3.3)

    If cubic and higher order terms in (t - t0 ) are neglected and the integral in t taken from - ∞ to + ∞, the result is

    (3.4)

    These operations can be justified* for sufficiently well behaved functions, ƒ.

    What happens when ƒ″(t0) also vanishes is itself an interesting question and in the context of path integration leads to an examination of caustics in electron optics. This is dealt with later in the book.

    Concerning the cubic, quartic, and higher powers of (t-t0), it is possible to give some heuristic arguments justifying their neglect. Consider the integral

    (3.5)

    We claim that λt³ and λt⁴ are small compared to 1 and will use a self consistency argument to demonstrate the point. Supposing they are small, K(λ) can be written

    (3.6)

    *Dropping regions of t in which ƒ′(t)≠0 can be justified as follows. If ƒ′(t) α≤ t β, then we can make the change of variables z =ƒ(t). Thus

    where

    Assuming that ϕ is differentiable, do an integration by parts to yield

    Hence F∝β goes to zero like 1/λ; from (3.4) regions having ƒ′=0 give contributions that decrease like 1/√λ and therefore ultimately dominate.

    where we have used integration formulas from the appendix to this section. Thus the higher powers of t yield terms that go to zero relative to the first term as λ→∞. Moreover, each factor becomes, after integration, 1/λ. In a sense this means that the effective size of t—or range of t contributing to the integral—is 1/√λ . With this way of counting, λt² is of order 1, λt³ of order λ-1/2, and so forth.c

    To summarize, as λ→∞, the behavior of the integral is determined solely by the point t0 where ƒ′(tmultiplying the function S[xfor which δS/δxis found to satisfy the Euler-Lagrange equations.

    →0 limit, what is done with it is not the same as is done in classical mechanics. Specifically it is used in a phase factor, so that if there are two classical paths converging on the same region we can still get interference patterns, in contrast to the purely classical prediction.

    Of course, to see interference patterns there must be something . Specifically, suppose there are two classical paths converging on each of the points ξ in some region. Let the associated classical actions for teach ξ be denoted S1(ξ) and S2(ξ). The propagator then has the general form

    (3.7)

    (actually these terms have real coefficients, which may not be equal). Since the probability distribution is given by |G|², an interference pattern arises from the effect of the cross term in the product GGc. The cross term is the real part of

    ei[S2

    Therefore, if the variation in S2 - S, the pattern will be destroyed by the rapid oscillation of this factor.

    Exercise: Suppose a wave function ψ has expectation values for position and momentum x1 and p1, and is sharply concentrated near these values, for example, in a minimum uncertainty wave packet. Let the system evolve under the Hamiltonian H=p²/2m+ V from time t1, to time t→0, at time t2 the wave function will be concentrated at x2, p2 where these are the values of position and momentum the classical system would have reached at t2, starting at x1,p1 at t1.

    Solution (very heuristic): Take

    (3.8)

    What’s needed is the propagator for a time interval T=t2-t1 and we guess (an educated guess) by analogy with (3.4) that it will take the form

    (3.9)

    with Sc for which δS=0 and with the appropriate boundary conditions. Moreover, C should be related to the second derivative of S but its exact form is not needed for the purpose of this exercise. What must be shown is that the amplitude

    is zero except when x2, p2 is the position in phase space at time t=t1+T of a particle leaving x1,p1 at t=t1. Then from (3.8) and (3.9)

    For uncertainties in x and p ). For the major contribution to the integral we look to the regions where ∂/∂ξ1, and ∂/∂ξ2 of the argument of the exponent vanish. Keeping ξ1, and ξ2 real, there are the separate requirements that

    ξ2=x2 ξ1=x1

    and

    (3.10)

    This says that (P2, x2) and (p1, x1) are related to each other by the canonical transformation generated by the function Sc. But Sc, generates (classically) time evolution, and we have obtained the condition that the particle stick to the classical trajectory. Just how closely it sticks is also implicit in these integrals.

    Closely related to the method of stationary phase is Laplace’s method, which deals with integrals of the form

    (3.11)

    with λ and g real. For g smooth and bounded from below points where g′(x)=0, g″(x)>0 give the greatest contribution. Because of the smallness of the exponential for large λ and g away from its minimum it is easier to prove asymptotic properties for G of (3.11) than for the F of (3.2) as λ→∞. Laplace’s method is dealt with in greater detail in Section 11.

    APPENDIX: USEFUL INTEGRALS

    (3.12)

    (3.13)

    (3.14)

    (3.15)

    (3.16)

    (3.17)

    (3.18)

    For Re a not positive these formulas can be interpreted by analytic continuation, the only complication being a branch point at a= 0.

    NOTES

    A standard source on asymptotic approximations is

    A. Erdelyi, Asymptotic Expansions, Dover, New York, 1956

    Additional references are given later in the book when we take up the subject of stationary phase and other asymptotic approximations in a more systematic way. At that stage we also discuss the extent to which the arguments given here can be rigorously applied to path integrals. The problem posed in the exercise of this section has been rigorously solved using methods closely related to the path integral, namely the Trotter product formula. See

    →0 Limit for Coherent States, Commun. Math. Phys. 71, 77 (1980)

    FOUR

    Vector Potentials and Another Proof of the Path Integral Formula

    In this section we extend the path integral to situations where a magnetic field is present. The extension has a surprising subtlety and gives us our first inkling of the intimate relation of path integrals to the theory of Brownian motion.

    First a generalization to three dimensions is needed. Nothing more strenuous is demanded for this generalization than making boldface all the position variables in (1.24). Each integral dxj is three dimensional and correspondingly the normalization factor becomes (m/2πi εN/2.

    For convenience in manipulation, and as a concession to popular usage, we shall also rewrite the formal path integral of (1.29) as follows

    (4.1)

    In the absence of magnetic fields we have already ascertained that S has the form

    (4.2)

    It will come as no surprise that in the presence of a magnetic field B, derivable from a vector potential A (B=∇ × A), the only change in (4.1) and (4.2) is that L becomes

    (4.3)

    In the action there will appear the additional term

    (4.4)

    which in any sensible theory should equal

    (4.5)

    The way to check whether the proposed propagator with the Lagrangian (4.3) is correct is to look at the expression for G as a limit

    (4.6)

    (In (4.6) N intermediate x’s are taken, rather than N — 1 as in (1.24).) From (4.5) it follows that the vector potential contribution is

    (4.7)

    but it is not clear whether A(x) should be evaluated at Xj, Xj. In fact, A We now prove that with the midpoint choice path integral theory gives the same result as Schrödinger’s equation, and later come back to show what would have happened had the wrong choice been made. In passing it will become evident why in (1.24) no fuss was made about the fact that xj, rather than xj+1 appeared as the argument of V(x).

    The way we check the validity of this method for handling magnetic fields is the same way Feynman first verified the entire formalism. Namely, take a wave function at time T, ψ(y, T), propagate it to time T+ε by means of the putative path integral propagator, and see if this evolution is the same as that given by the Schrödinger equation. For this purpose we take as propagator the quantity under the limit operation in (4.6) for the case Ncontribution. This will be used to propagate ψ for a time ε and we shall eventually let ε go to zero to allow comparison with the ∂/∂t operation of Schrödinger’s equation. If the theory is correct we should therefore have

    (4.8)

    The next few steps involve a small nightmare of Taylor expansions and Gaussian integrals. Define the variable

    ξ=y—x

    Let this be the variable of integration in (4.8) and expand ψ(y, T), A, and V about ξ=0. Equation 4.8 becomes

    (4.9)

    t~1/√λ. To be precise, we are considering the integral (4.9) in the limit ε→0 and are keeping both the lowest terms (which turn out to be of order unity and cancel exactly) and the first term in ε (which will be of order ε), but no higher powers. The argument of the exponent, imξ²/2 ε, will be of order unity which is to say ξ is of order √ε. This is demonstrated exactly as was done in the previous section, namely by doing Gaussian integrals.

    First we note that the term exp(-iεV) can be factored from the integral—it does not depend on ξ—and can also be written as (1 - iεV) to first order in ε.

    Next the factor exp(-;∇V·ξ) can simply be dropped. When the integration overt ξ is performed this will be at most of order ε³/². This, by the way, is the reason that it does not matter whether xj or xj+1 is taken as the argument of V.

    The other factors all remain, and we expand the exponential involving A up to order ξ², that is, up to ε. Equation 4.9 becomes

    (4.10)

    It is understood in (4.10) that repeated indices k, l, m, n are summed from 1 to 3. By symmetry in ξ, all integrals over odd powers of ξ as well as all products of the form ξnξl, nl, give zero upon integration. Furthermore, each ξnξn (no summation) gives

    (4.11)

    by the appendix to Section 3 (3.13). The only cross term in (4.10) is thus from ψ and integration yields

    (4.12)

    The cancellation of the normalization factor is of course no accident. ψ(x, T) is brought to the left hand side, which becomes

    (4.13)

    We divide the entire equation by ε i. The result is

    (4.14)

    which is the same as

    (4.15)

    with p the operator - i . Therefore wave functions propagated by the path integral have a time evolution satisfying Schrödinger’s equation.

    NOTES

    This section follows Feynman’s original paper [Rev. Mod. Phys. 20, 367 (1948)] quite closely.

    FIVE

    The Ito Integral and Gauge Transformations

    Let us examine the consequences of an incorrect evaluation of

    (5.1)

    Suppose this were interpreted as

    (5.2)

    by (ξ·∇)A. Following the consequences of this replacement to (4.12), we find (εe/mc∇·A instead of half that quantity. Finally this will have the effect of adding an additional term

    (5.3)

    to the Schrödinger equation, a term that does not belong there.

    You should receive these remarks with the greatest alarm. The basic fact about Riemann integrals, the fact that makes them well defined objects, is the property that approximating sums

    (5.4)

    , so long as maxj(tj+1— tj)→0, no matter what points tj* are selected (for tj+1≥tj*≥tj, tN+1 = b, t0 =a). It follows that whatever sense we ultimately make out of the integral (5.1), it will not be a Riemann integral.

    In the study of Brownian motion, K. Ito found similar problems in attempting to define an integral

    (5.5)

    where x(t) is a Brownian motion path, that is, a particular trajectory followed by a particle undergoing Brownian motion. He chose to define this integral as

    (5.6)

    He then found the remarkable result

    (5.7)

    x=x(tt)-x(txt t (in mathematics—not in physical Brownian motion). The velocity, or derivative of xxtt t→0. In this light let us examine (5.7). Consider the identity

    (5.8)

    with the same conventions as before for xj, and so on. ϕ is assumed to be a well behaved function and (5.8) can be rewritten

    (5.9)

    where uθj is defined as

    (5.10)

    and 1≥θ≥0. Suppose we let θ = 0; then the first sum in (5.9) becomes, in the limit n→∞, the Ito integral. But the important thing to notice is that as n→∞ xt x)² is ϕ″ dt the midpoint rule—the anomalous term would not appear.

    In quantum mechanics we have found the need for the midpoint rule and by now the reason for our difficulties should be evident. We found in Section 4 that the size of ξ is √ε. Thus the broken line paths that we

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