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PHYSICS
CSIR NET

THIRD EDITION
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 Physics Books available for NET Preparation

Mathematical Physics.
1. Mathematical Methods in the Physical Sciences : Mary L. Boas link
2. Advanced Engineering Mathematics: Erwin Kreyszig link
3. Mathematical Physics: H. K. Dass. link

Classical Mechanics
1. Classical Mechanics : J.C. Upadhyaya link
2. Classical Mechanics : Herbert Goldstein link
3. An Introduction to Mechanics: Kleppner & Kolenkow link
4. Concept of Physics (Volume I): H. C. Verma link

Electromagnetic Theory
1. Introduction to Electrodynamics: David J. Griffiths link
2. Classical Electrodynamics: Walter Griener link

Quantum Mechanics
1. Quantum Mechanics Concepts & Applications: Nouredine Zettili link
2. Introduction to Quantum Mechanics: David J. Griffiths link
3. Quantum Physics: H.C. Verma. link
4. Quantum mechanics: 500 problems with solutions: G. Aruldhas link

Thermodynamics and Statistical Physics

1. Fundamentals of Statistical & Thermal Physics: F. Rief. link
2. Statistical Mechanics: R. K. Patharia link
3. Thermal Physics: Garg, Bansal, Ghosh. link

Electronics & Experimental Methods

1. Digital Electronics: Malvino & Leach link
2. Electronic Devices & Circuits: Jacob Millman & Christos C. Halkias link
3. Electronic Devices and Circuit Theory 7tn Robert L Boylestad, Louis Nashelsky Link
Atomic & Molecular Physics
1.Introduction to Atomic Spectra: Harvey Elliott White. link
2. Fundamentals of Molecular Spectroscopy Colin N.Banwell, Elaine M.MaCash link
3. Atomic& Molecular Spectra: Laser by Rajkumar Link

Condensed Matter Physics

1. Introduction to Solid State Physics: Charles Kittel. link
2. Solid State Physics: Puri & Babbar. link
3. Solid State Physics: M. A. Wahab link

Nuclear and Particle Physics

1. Nuclear Physics, An Introduction: S. B. Patel. link
2. Introduction to Elementary Particles: David J. Griffith. link

For preparation please use the following handwritten notes

and materials
For those who start NET preparation. Easiest and organized way to qualify is to go through these
materials and prepare.
➢ Handwritten materials Link +All the associate files to learn fundamental.
➢ And Materials {To learn theory in detail}
➢ To do problems using assignment sets
➢ Assignment Career Endeavour
➢ Assignment FIZIKS
➢ FINALLY subject wise previous year questions
➢ Previous Year question

Best of luck. Hard work will pay off👍

 Syllabus for NET
CSIR-UGC National Eligibility Test (NET) for Junior Research
Fellowship and Lecturer-ship PHYSICAL SCIENCES

PART ‘A’ CORE

I. Mathematical Methods of Physics

Dimensional analysis. Vector algebra and vector calculus. Linear algebra, matrices, Cayley-Hamilton
Theorem. Eigenvalues and eigenvectors. Linear ordinary differential equations of first & second
order, Special functions (Hermite, Bessel, Laguerre and Legendre functions). Fourier series, Fourier
and Laplace transforms. Elements of complex analysis, analytic functions; Taylor & Laurent series;
poles, residues and evaluation of integrals. Elementary probability theory, random variables,
binomial, Poisson and normal distributions. Central limit theorem.

II. Classical Mechanics

Newton’s laws. Dynamical systems, Phase space dynamics, stability analysis. Central force motions.
Two body Collisions - scattering in laboratory and Centre of mass frames. Rigid body dynamics-
moment of inertia tensor. Non-inertial frames and pseudoforces. Variational principle. Generalized
coordinates. Lagrangian and Hamiltonian formalism and equations of motion. Conservation laws and
cyclic coordinates. Periodic motion: small oscillations, normal modes. Special theory of relativity-
Lorentz transformations, relativistic kinematics and mass–energy equivalence.

III. Electromagnetic Theory

Electrostatics: Gauss’s law and its applications, Laplace and Poisson equations, boundary value
problems. Magnetostatics: Biot-Savart law, Ampere's theorem. Electromagnetic induction. Maxwell's
equations in free space and linear isotropic media; boundary conditions on the fields at interfaces.
Scalar and vector potentials, gauge invariance. Electromagnetic waves in free space. Dielectrics and
conductors. Reflection and refraction, polarization, Fresnel’s law, interference, coherence, and
diffraction. Dynamics of charged particles in static and uniform electromagnetic fields.

IV. Quantum Mechanics

Wave-particle duality. Schrödinger equation (time-dependent and time-independent). Eigenvalue

problems (particle in a box, harmonic oscillator, etc.). Tunneling through a barrier. Wave-function in
coordinate and momentum representations. Commutators and Heisenberg uncertainty principle.
Dirac notation for state vectors. Motion in a central potential: orbital angular momentum, angular
momentum algebra, spin, addition of angular momenta; Hydrogen atom. Stern-Gerlach experiment.
Time-independent perturbation theory and applications. Variational method. Time dependent
perturbation theory and Fermi's golden rule, selection rules. Identical particles, Pauli exclusion
principle, spin-statistics connection.

V. Thermodynamic and Statistical Physics

Laws of thermodynamics and their consequences. Thermodynamic potentials, Maxwell relations,

chemical potential, phase equilibria. Phase space, micro- and macro-states. Micro-canonical,
canonical and grand-canonical ensembles and partition functions. Free energy and its connection with
thermodynamic quantities. Classical and quantum statistics. Ideal Bose and Fermi gases. Principle of
detailed balance. Blackbody radiation and Planck's distribution law.

VI. Electronics and Experimental Methods

Semiconductor devices (diodes, junctions, transistors, field effect devices, homo- and hetero-junction
devices), device structure, device characteristics, frequency dependence and applications. Opto-
electronic devices (solar cells, photo-detectors, LEDs). Operational amplifiers and their applications.
Digital techniques and applications (registers, counters, comparators and similar circuits). A/D and
D/A converters. Microprocessor and microcontroller basics. Data interpretation and analysis.
Precision and accuracy. Error analysis, propagation of errors. Least squares fitting,
 PART ‘B’ ADVANCED

I. Mathematical Methods of Physics

Green’s function. Partial differential equations (Laplace, wave and heat equations in two and three
dimensions). Elements of computational techniques: root of functions, interpolation, extrapolation,
integration by trapezoid and Simpson’s rule, Solution of first order differential equation using Runge-
Kutta method. Finite difference methods. Tensors. Introductory group theory: SU(2), O(3).

II. Classical Mechanics

Dynamical systems, Phase space dynamics, stability analysis. Poisson brackets and canonical
transformations. Symmetry, invariance and Noether’s theorem. Hamilton-Jacobi theory.

III. Electromagnetic Theory

Dispersion relations in plasma. Lorentz invariance of Maxwell’s equation. Transmission lines and
wave guides. Radiation- from moving charges and dipoles and retarded potentials.

IV. Quantum Mechanics

Spin-orbit coupling, fine structure. WKB approximation. Elementary theory of scattering: phase
shifts, partial waves, Born approximation. Relativistic quantum mechanics: Klein-Gordon and Dirac
equations. Semi-classical theory of radiation.

V. Thermodynamic and Statistical Physics

First- and second-order phase transitions. Diamagnetism, paramagnetism, and ferromagnetism. Ising
model. Bose-Einstein condensation. Diffusion equation. Random walk and Brownian motion.
Introduction to nonequilibrium processes.

VI. Electronics and Experimental Methods

Linear and nonlinear curve fitting, chi-square test. Transducers (temperature, pressure/vacuum,
magnetic fields, vibration, optical, and particle detectors). Measurement and control. Signal
conditioning and recovery. Impedance matching, amplification (Op-amp based, instrumentation amp,
feedback), filtering and noise reduction, shielding and grounding. Fourier transforms, lock-in
detector, box-car integrator, modulation techniques. High frequency devices (including generators
and detectors).

VII. Atomic & Molecular Physics

Quantum states of an electron in an atom. Electron spin. Spectrum of helium and alkali atom.
Relativistic corrections for energy levels of hydrogen atom, hyperfine structure and isotopic shift,
width of spectrum lines, LS & JJ couplings. Zeeman, Paschen-Bach & Stark effects. Electron spin
resonance. Nuclear magnetic resonance, chemical shift. Frank-Condon principle. Born-Oppenheimer
approximation. Electronic, rotational, vibrational and Raman spectra of diatomic molecules, selection
rules. Lasers: spontaneous and stimulated emission, Einstein A & B coefficients. Optical pumping,
population inversion, rate equation. Modes of resonators and coherence length.

VIII. Condensed Matter Physics

Bravais lattices. Reciprocal lattice. Diffraction and the structure factor. Bonding of solids. Elastic
properties, phonons, lattice specific heat. Free electron theory and electronic specific heat. Response
and relaxation phenomena. Drude model of electrical and thermal conductivity. Hall effect and
thermoelectric power. Electron motion in a periodic potential, band theory of solids: metals,
insulators and semiconductors. Superconductivity: type-I and type-II superconductors. Josephson
junctions. Superfluidity. Defects and dislocations. Ordered phases of matter: translational and
orientational order, kinds of liquid crystalline order. Quasi crystals.

IX. Nuclear and Particle Physics

Basic nuclear properties: size, shape and charge distribution, spin and parity. Binding energy, semi-
empirical mass formula, liquid drop model. Nature of the nuclear force, form of nucleon-nucleon
potential, charge-independence and charge-symmetry of nuclear forces. Deuteron problem. Evidence
of shell structure, single-particle shell model, its validity and limitations. Rotational spectra.
Elementary ideas of alpha, beta and gamma decays and their selection rules. Fission and fusion.
Nuclear reactions, reaction mechanism, compound nuclei and direct reactions. Classification of
fundamental forces. Elementary particles and their quantum numbers (charge, spin, parity, isospin,
strangeness, etc.). Gellmann-Nishijima formula. Quark model, baryons and mesons. C, P, and T
invariance. Application of symmetry arguments to particle reactions. Parity non-conservation in weak
interaction. Relativistic kinematics.
ESSENTIAL
NOTES
PHYSICS
CSIR NET

THIRD EDITION
Follow Us on Telegram and
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Those who prepare for NET Exam. Please go

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 PHYSICS
Part – BCSIR
& C NET
Unit-1

BASIC MATHEMATICAL METHODS

Vector Algebra Applying these products to equations (1), (2)
Vectors only in three mutually perpendicular and (3);
Cartesian co-ordinate system can be shown as → →
A · B = a1 b1 + a2 b2 + a3 b3 ;
→
A = a1^i + a2^j + a3 ^k …(1) → →
→ ^ ^ ^ B · C = b1 c1 + b2 c2 + b3 c3 etc.
B = b1 i + b2 j + b3 k …(2)
Also
→ ^ ^ ^
C = c1 i + c2 j + c3 k …(3) ⎪⎪ ^i ^j ^k ⎪⎪
→ → → ⎪a ⎪
^ ^
D =d i +d j +d k ^ …(4) (A × B ) = ⎪ 1 a2 a3 ⎪
1 2 3 ⎪⎪ b1 b2 b3 ⎪⎪
where ^i , ^j , k^ are the unit vectors in the directions
x, y and z respectively. ^ ^
= i (a2 b3 – a3 b2 ) + j (a3 b1 – a1 b3 )
The vector multiplication is executed in two
ways— + k^ (a b – b a )
1 2 1 2

(a) Scalar or dot product : This product which is a determinant.

supplied the scalar resultant. Scalar product is
→→ → → → → Vector Triple Product
commutative, i.e., A · B = B · A , A · B = A · B
cos θ, where θ is the angle between the vectors. In → → →
By expanding the vectors A , B and C in
→ →
other words A · B is equal to the multiplication of terms of ^i , ^j , ^k and using the results of unit vector
one vector with the projection of other over it. products, we can obtain,
Also ^i · ^j = ^j · ^i = ^i · k^ = ^k· ^i = ^j · ^k = ^k· ^j = 0 → → → → → →
A ( B · C ) ≠ (A · B ) C
^ ^ ^ ^ ^ ^
and i · i = j · j = k· k = 1 → → →
(b) Vector or cross product : As its name and A · (B × C ) represents the volume of paral-
→ →
suggests cross product results in vector output. lelopiped having the base arms B & C and height
This product is not commutative →
A.
→ → → → → →
i.e., A × B = – ( B × A ) ≠ ( B × A ). → → → → → → → → →
In the cross product the order of multiplication A · (B × C ) = B · ( C × A ) = C · ( A × B )
plays a significant role. → →
are the scalar triple product of vectors A , B
→ → →
A × B = A · B · sin θ · ^n and C can be found as
where θ is the angle between the vectors A and B.
→ → → → → ⎪⎪ a1 a2 a3 ⎪⎪
And n^ is the normal unit vector. ( A × B ) vector A · ( B × C ) = ⎪ b1 b2 b3 ⎪
⎪⎪ ⎪⎪
→ (Scalar or box c1 c2 c3
lies perpendicular the plane containing vector A product)
→
and B . = Volume of parallelopiped.
66 | Physical Sciences

→ → → → → → Divergence of a Vector Field

whereas A × ( B × C ) ≠ ( A × B ) × C is called a
vector triple product. In general ^ ^ ^
If V = Vx i + V y j + V zk represents a vector
→ → → → → → → → → field, then divergence of V can be given as
A × (B × C ) = ( A · C ) B – ( A · B ) C div V = ∇ · V
(this is a vector) ∂ ∂ ∂
→ → → →→ → → → →
(
= ^i + ^j + ^k
∂x ∂y ∂z )· (Vx ^i + Vy j^ + Vz k)
^
(A × B ) × C = (A · C ) B – (B · C ) A
∂Vx ∂Vy ∂Vz
(this is a vector) = ( ∂x
+
∂y
+
∂z )
→ → → → → → (A scalar quantity)
A · (B × C ) = (B × C ) · A
If the divergence of any vector field vanishes
(since dot product is commutative) then the field is termed as solenoidal field.
→ → → → → → Curl of a Vector Field
But A · ( B × C ) = B · ( C × A )
(property of box product.) If V = Vx ^i + Vy^j + Vz ^k representes a vector
field, then the curl of V can be given by
→ → → → → → Curl V = ∇ × V
Hence ( B × C ) · A = B · ( C × A )
^i ^j ^k 
This means that the dot and cross can be
∂ ∂ ∂ 
 ∂x 
interchanged in a box product or scalar triple =
∂y ∂z
V 
product.
x Vy Vz
Gradient, divergence and curl of scalar
and vector functions— Curl of a vector field is a vector and if the
curl of any vector vanished the field is known as
The ∇ (known as del or delta) is a three irrotational field.
dimensional operator,
Vector Integral Theorems
∂ ^ ∂ ^ ∂
∇ ≡ ^i +j +k The divergence theorem of Gauss, the Stoke’s
∂x ∂y ∂z theorem and Green’s theorem and related integral
which can be operated on a scalar function or a theorems :
vector function. Del shows the rate of change of a (a) Gauss’ divergence theorem—It is an
scalar or vector field w.r.t. three mutually perpen- important theorem to relate volume integral of a
dicular axis. vector to surface integral of this vector.
According to this theorem the volume integral
The Gradient of divergence of any vector is equal to the surface
integral of that vector for the surface enclosing the
If a scalar field is represented by φ (x, y, z), volume i.e.,
the rate of change of φ with the distance s is given → → → →
by ∫∫∫ V
( ∇ · A ) dV = ∫∫ A · dS
S
dφ ∂φ dx ∂φ dy ∂φ dz (b) Stoke’s theorem—According to Stoke’s
= · + · + ·
ds ∂x ds ∂y ds ∂z ds theorem the line integral of a vector along a closed
path C is equal to the surface integral of the curl
= (^i ∂x∂φ + ^j ∂y∂φ + ^k ∂z∂φ)· (^i dsdx + ^j dsdy + ^k dsdz) of that vector taken, over the surface or

∫
→ → → →
The quantity in first bracket is known as
O A · dr =
C
∫∫ ( ∇ × A ) · ^n dS
S
Gradient of φ and is written as ∇φ or grad φ. →
A has continuous derivatives, it is a vector
Again, we can say
function of position. The surface S is open two
∇φ = (^i ∂x∂ + ^j ∂y∂ + ^k ∂z∂ ) φ sided. C is a simple non-intersecting curve. C is
traversed in the positive direction.
 Physical Sciences | 67

→ Line integral, surface integral and volume

(c) Green’s theorem—Let the vector field A integral have significant importance in vector
is the product of a scalar function φ 1 and the
integration.
gradient of another scalar function φ2 i.e.,
A = φ1 ∇ φ2
L.I. =
∫P2 A · dr
P1

∫P2 ^ ^ ^ ^ ^ ^
then according to Green’s theorem; = (a1 i + a2 j + a3k) (dx i + dy j + dzk)
∫∫∫
P1
(φ2 ∇2 φ1 + ∇φ2 · ∇φ1) dV
V = ∫P2
P1
(a1 dx + a2 dy + a3 dz)
= ∫∫ S
(φ 2 ∇ φ1) dS = ∫P2
P1
a1 dx + ∫P2a dy
P1
2 + ∫
P2
P1
a3 dz
This is the first form of Green’s theorem P 1 and P2 are points along the path or line.
and ∫∫∫ V
(φ1 ∇2 φ2 – φ2 · ∇2 φ1) dV Matrices
= ∫∫ S
(φ 1 ∇ φ2 – φ2 ∇ φ1) dS A rectangular array of m, n numbers (m being
the number of rows and n columns) is called a
This is second part of Green’s theorem. matrix of order m × n. e.g.,

Vector Differentiation [ 5 –1 6
]
2 0 3 is a 2 × 3 matrix.
The rules for vector differentiations are no The elements of a matrix are represented as
more different from those for scalar differentiation. [aij] or (aij). Where i and j represent its row and
But here a change in the physical parameter column position.
occurs not only a change in the magnitude but also  aa11 aa12 … a1n 
a change in the direction or both.
A=  21 22 … a2n 
 
....

→
e.g., r , the position vector of a point is given  am1 am2 … amn  is a m × n matrix.
by
Special Matrices
→
r (t) = ^i x (t) + ^j y (t) + ^kz (t) (a) A matrix having only one row is known as
→ raw matrix or vector. Likewise a column matrix
dr dx (t) ^ dy (t) ^ dz (t)
and = ^i +j +k have a single column.
dt dt dt dt (b) A complex matrix may contain real as well
→ as complex elements.
V = ^i Vx + ^j Vy + ^kVz (c) Square matrix are those for which number
which will be velocity. of rows are exactly equal to the number of column
d dB dA i.e., m = n otherwise the matrix is known as
Also (A × B) = A × + ×B rectangular matrix.
dt dt dt
(d) Null matrix contains all elements equal to
where A and B are vectors zero.
d dφ dA (e) A square matrix which have every element
(φA) = A+ φ (φ is a scalar) of the main diagonal equal to unity and all other
dt dt dt
element equal to zero is, termed as diagonal matrix
Vector Integration amn = 0 for m ≠ n.
1 0 0 0 
The rules of vector integration are also similar 0 
 1 0 0 
to those for scalar integration.
0 0 1 0 
r (t) = ^i x (t) + ^j y(t) + ^kz(t)

If 0 0 0 1  is a 4 × 4 diagonal matrix.
Then (f) A square matrix whose determinant is zero
∫r(t) dt = ^i ∫x(t) dt + ^j ∫y(t) dt + k^ ∫z (t) dt is called singular; all non-square matrices are also
called singular.
68 | Physical Sciences

Another collection of special matrices is listed Cramer’s Rule

below :
It is an easy way to find the solution of n
Name of Notation How to get it from equations having n variables.

    
Matrix for it original Matrix A

a11 a21 … a1n x1 r1
Transpose A′ or At Interchange rows and

a21 a22 … a2n

x2 r2   
of A column in A a31 a32 … a3n x3 = r3   
–
    

....

....

....
a  x
A or A*
 r 
Complex Take the complex conju-
Conjugate gate of each element n1 an 2 … ann n n
of A
–
A X R
Transpose A′ or A† Take complex conjugate or AX = R
Conjugate (A dagger) of each element and trans- ⇒ A–1AX = A–1R
of A pose it.
⇒ X = A–1R
Adjoint or adj A Replace all elements by
Adjugant their corresponding co-
To solve for x1, x 2 , x 3 …xn in terms of r1 , r 2 ,
of A factors r3 …rn and the known matrix A.
Reciprocal A–1 Divide each element of a11 a21 … a1n
or Inverse adj A by det A a21 a22 … a2n
of A adj A Let ∆ = a31 a32 … a3n
A–1 =

....
....
|A|
an 1 an 2 … ann
If a square matrix A is same as its transpose = Determinant of A
conjugate then A is Hermitian matrix, i.e., A = A†
also if A = A– † then the matrix is Skew-Hermitian. Then according to Cramer’s rule, we have
∆ ∆ ∆
Principal diagonal of the matrix is the diago- x1 = 1 ; x2 = 2 ; … xn = n
nal containing the elements a11, a22, a33…ann. The ∆ ∆ ∆
n where ∆ k is the determinant, obtained from ∆ by
sum of all these elements i.e., Σ all is known as replacing the kth column by the column vector R.
l=1
the trace of matrix. Eigen Values and Eigen Vectors of
A matrix is known to be singular if its Operators
determinant is zero, i.e., if det (A) ≠ 0, A is a non- Matrix are used as operators. Maxtrix operated
singular matrix. on a column vector X, produces a vector in the
A complex matrix is unitary if its complex same direction, but may be of a different
conjugate transpose is equal to its inverse. magnitude.
A† = A–1 or AA–1 = I Thus, AX = λX
A real matrix (A) is said to be orthogonal if The operation of A on X is only to change its
magnitude by λ times.
At = A–1 , AA–1 = I
The λ is the eigen value of operator A,
A real unitary matrix is an orthogonal matrix.
belonging to the eigen vector X.
Rank of a Matrix (a11 – λ) x 1 + a12 x2 + a13x3 + … + a1nxn = 0
The rank of a matrix is the order of the largest a21x1 + (a22 – λ) x 2 + a23x3 + … + a2nxn = 0
non-vanishing determinant. For A n × n matrix we
first find the determinant for order n if it is zero, ……………………………………………
then for order n – 1 and so on. If all the deter- ……………………………………………
minant for order (n – 1) are zero then evaluate
……………………………………………
determinant for (n – 2) order. If even one non-zero
determinant is there the rank will be (n – 2). an 1 x1 + an 2 x2 + an 3 x3 + … + (ann – λ) x n = 0
 Physical Sciences | 69

These equations have non-trivial solutions, if that the co-ordinate (u, v) of point P, referred to
and only if new axes are related to (x, y) by
u = x cos θ + y sin θ
a a – λ a12 … a 1n 
 · =0 v = – x sin θ + y cos θ
11
a22 – λ … a 2n
 · 
21
·· A transformation from one set of rectangular
a n1 an 2 … ann – λ  co-ordinates to another set of rectangular co-
ordinate is called an orthogonal transforamation.
This is a polynomial of degree n in λ. The An important property of this transformation
roots of this equation are called eigen values or is that
characteristic values of matrix, A. For each eigen u2 + v2 = x2 + y2
value a non-trivial solution X ≠ 0, X is called a
eigen vector. A general transformation that will be obtained
when the new axes are rotated through different
Characteristics of Eigen Value angles θ and φ and the angle between them does
1. Any matrix A and its transpose both have not remain a right angle.
the same eigen value. Another type of transformation is
2. The trace of the matrix equals to the sum
of the eigen values of a matrix. u = ax
3. The determinant of the matrix A equals v = by
to the product of the eigen values of A More general form of linear transformation in
4. If λ 1 , λ2, … λ n are n-eigen values of A, two dimension is
then u = a1 x + b1 y
(i) Eigen values of k A are k λ1, k λ2 , … k λn
v = a2 x + b2 y
1 1 1
(ii) Eigen values of A– 1 are , , …
λ1 λ2
(iii) Eigen values of A are λ1 , λ2 , … λn
m m m
λn
m
[] [
u
v = a2 b2
a1 b1
][ ] x
y
5. If λ is an eigen value of a matrix A, then Similarly
1 Y = AX
is the eigen value of A– 1.
λ
 
y1
Note—  
y2
● If λ be an eigen value of a non-singular Y =
 
·· ;
matrix, then
|A|
λ
is an eigen value of Adj  
yn
A.  a11 a12 … a1n 
● If λ be an eigen value of a non-singular  a21 a22 … a 2n 
matrix, then λ – 1 is the eigen value of A =
 ·· ··  ,
A– 1 .
● Eigen values of a Hermitian matrix are
 an 1 an 2 … 
ann

real.  
x1
● Characteristics root of a diagonal matrix  
x2
are diagonal elements of the matrix. X =
 ···
● Characteristics root of a upper triangular
or lower triangular matrix is its diagonal
 
xn
element. |A| is determinant of transformation matrix A,
is called the modulus of the transformation.
Linear Transformations
Let P(x, y) be a point in xy plane. If we take |A| = 0, transforamation is singular
another set of rectangular axis inclined to the |A| ≠ 1, then transformation is non-singular or
former set at an angle θ, then it is easy to show regular.
70 | Physical Sciences

Orthogonal Transformations Hence, we have

Any square matrix A is called orthogonal
matrix if AA′ = I, Y = AX is said to be orthogonal
dy
dx e
∫ P(x) dx + P(x) y e∫ P(x) dx = Q(x) e∫ P(x) dx
if A is in orthogonal matrix.
The determinant of orthogonal matrix is + 1
or
d
[
dx y e
∫ P(x) dx = Q(x) e∫ P(x) dx
]
or – 1.
or y = e–∫ P(x) dx [∫Q(x) e∫ P(x) dx + C]
| A′A | = | I | i.e., | A |2 = 1
|A| = ± 1 Perfect Differential
| A | = + 1 proper Equation Mdx + Ndy = 0 is said to be perfect
| A | = – 1 Improper ∂M ∂N
differential if = · This equation for which
∂y ∂x
Linear Differential Equations ∂M ∂N
≠ is not perfect differential.
A differential equation written in the form ∂y ∂x
dny dn – 1y dn – 2y The solution of this exact equation is given by
a0 (x) n + a1 (x) n – 1 + a2 (x) n – 2 + …
dx dx
dy
dx
∫ M dx + ∫ N dy = C
+ an – 1 + an (x) y = R (x) y = constt. only those term which
dx are independent of x
is known as linear differential equation. In a linear Some important partial differential equations :
differential equation all the derivatives of first
power, they should not be quadratic, cubic or (a) Wave equation—
higher powered. 1 ∂2 x
∇2 x = 2 2
d2 y dy ν ∂t
e.g., 2x 2 + 6 – 6xy = cos x 1 ∂2
dx dx
is a second order linear equation, whereas
or ( )
∇2 – 2 2 x = 0
ν ∂t
2
d2 y dy 3 or x = 0
y 2–x
dx () dx
+ x2y = e –2x 1 ∂ 2
where 2 ≡ ∇2 – 2 2 is D’Alembertian operator,
is a third order non-linear equation. ν ∂t
which is a four dimensional analogue of the
In the form of differential operator D the Laplacian in Minkowski space.
linear differential equation can be written as
In case of electromagnetic waves x may be a
[a0(x) Dn + a1 (x) Dn – 1 + … + an – 1 D + an(x)]y component of E or H, also v is replaced by ‘c’, the
= R(x) speed of light in free space.
Solution of linear differential equation, (b) Laplace’s equation—
dy ∇2 u = 0
+ P(x)y = Q(x).
dx The function u may be the gravitational
[P(x) y – Q(x)] dx + dy = 0 = M dx + N dy potential in a region containing no-matter, the
electrostatic potential in a charge free region, the
where M = P(x) y – Q(x) steady state temperature or the velocity potential
and N = 1 for an incompressible fluid with no vertices, no
∂M sink or sources.
My = = P(x), Nx = 0
∂y (c) Poisson’s equation—
1 ∇2 u = f (x, y, z)
Hence, f(x) = P(x) = (My – Mx)
N It is a modification of Laplace’s equation, u
= [P(x) – 0] may be a gravitational potential containing matter
Hence, integrating factor or the electrostatic potential in a region of charge.
The function f (x, y, z) is called the source of
= e∫ P(x) dx density.
 Physical Sciences | 71

(d) The heat flow equation—The diffusion
or heat equation is of the form
1 ∂u
∇2 u =
α2 ∂t
u may be a non-steady state temperature is a
region with no heat source.
The quantity α2 is a constant known as
diffusivity.
(ii) Real and equal roots—If the both roots
are real and equal say m, then the general solution
be
y = (Ax + B) emx
(iii) Complex conjugate roots—Now if the
roots of auxilliary equation are α ± iβ, the general
solution be
y = eαx [A cos βx + B sin βx]

(e) Helmholtz equation— here A and B are arbitrary constants.

∇2 u ± k2u = 0 (b) Inhomogeneous equations—An example
of inhomogeneous equation is shown below :
where u represents the time independent part.
d2 y dy
+ a1 + a2 y = X
Bernoulli’s Equation dx2 dx
dy where a1 and a2 are constants and X is a function
The equation + Py = Qyn , where P and Q of x. Solution of such equation is a combination of
dx
are constants or functions of x and y alone and n is two parts. First, the complementry function which
an integer other than zero or unity, is called can be obtained for (X = 0) as in case of
extended form of linear equation or Bernoulli’s homogeneous equation and second the particular
equation. X
integral y = ·
f(D)
Solution of Second Order Linear For different value of x, P. I. will be given as :
Differential Equations with Constant Form of X Particular Integral
Coefficients
eax 1 ax
(a) Homogeneous equations—A homoge- e If F(a) ≠ 0
F(a)
neous equation is given by
sin ax 1 1
d2 y dy sin ax = sin ax
+ a1 + a2 y = 0 F(D2) F(– a2)
or cos ax
dx2 dx
If F(– a2 ) ≠ 0
where a1 and a2 are constants.
1 1
d cos ax = cos ax
In terms of operator D ≡ , the above equa- F(D2) F(– a2)
dx
tion may be written as If F(– a)2 ≠ 0
(D 2 + a1 D + a2 ) y = 0 (auxilliary equation) xm y = {f (D)}–1 xm, on expanding

or D 2 + a1 D + a2 = 0 {f (D)}–1 binomially in the

ascending powers of D up to with
or f (D) = 0 power and then operating with
each operator on xm.
which is a quadratic equation in D, roots of which
can be real and equal, real and unequal or a V e ax , V is 1 ax ax 1
complex pair. any function y = f(D) (V e ) = e f(D + a) (V)
of x
Now, we will discuss the case separately.
(i) Real and unequal roots—If the roots of
given equation are real and distinct say m1 and m2
xV
y= [ x–
1
f(D)
· f ′(D) ] 1
f (D)
(V)

the general solution of equation be The required solution of equation will be

y = C 1 em1x + C2em2x y = C. F. + P. I.
72 | Physical Sciences

Special Functions (i) Assume the solution to be

y = x m (a0 + a1 x + a2 x 2 + … an x n)
To Solve the Equation
dy d 2y
d 2y dy (ii) Then, substitute y, and in
P0 + P1 + P2 y = 0 …(1) dx dx 2
dx 2 dx equation (1)
where P’s are polynomials in x and P0 ≠ 0 at x = 0 (iii) Equate to zero the coeff. of the lowest
(i) Assume its solution to be of the form degree term in x. It gives a quadratic equation as
the indicial equation.
y = a0 + a1 x + a2 x 2 + … + a n x n 2 …(2)
(iv) Equating to zero the coeff. of the other powers
dy , d 2y of x, find the values of a1, a2 , a3 , … in terms of a0.
(ii) Calculate from equation (2) and
dx dx 2
substitute in (1) Solution about Singular Points
(iii) Equate to zero the coefficient of the There are other two types of singular points
various powers of x and determine a 2, a3 , a4 in (1) Regular singular points
terms of a0, a1 . (The result obtained by equating to (2) Irregular singular points
zero is, the coefficient of x n that is called the (I) y′′ + P1(x) y′ + P2(x)y = 0 …(5)
recurrence relation). and assume that at least one of the function P1 and
(iv) Substituting the value of a2 , a 3 , a 4 in P 2 is not analytic (P 1 = ∞ or P2 = ∞ ) at x = a, so
equation (2) and the desired series solution having that x = a is a singular point of (5)
a0 , a1 as its arbitrary constant. Consider Q1(x) = (x – a) P1(x)
Validity of series solution of the equation Q2(x) = (x – a)2 P0 (x)
d 2y dy If Q1 and Q 2 are analytic (not ∞) at x = a, then
P 0 (x) 2 + P1(x) + P2(x) y = 0 …(3)
dx dx x = a is called a regular singular point, otherwise
If P0 (a) ≠ a, then x = a is called an ordinary irregular.
point of equation (3) otherwise it is a singular Frobenious Method
point
A Singular point x = a of (3) is called regular If x = 0 is a regular singular point, then
if, when equation (3) is put in the form soltuion will be
∞
d 2 y Q1(x) dy Q2 (x) y = Σ ak xm + k
+ + =0 k=0
dx 2 x – a dx (x – a)2
dy , d 2y
where Q1 (x) and Q2 (x) posses derivative of all y, substitute and determine value of m.
dx dx 2
orders in the neighbourhood of a.
On equating the coefficient lower power of x
(i) When x = a is an ordinary point of (3), its
in the identity to zero, a quadratic equation is m
every solution can be expressed in the form
(indicial equation is obtained.)
y = a0 + a1 (x – a) + a2 (x – a)2 + … …(4)
(ii) When x = a is an regular singularity of Bessel’s Equation
(3), at least one of solution can be expressed as d 2y dy
x 2 2 + x + (x 2 – n 2 )y = 0
m 2
y = (x – a) [a0 + a1 (x – a) + a2 (x – a) + …] dx dx
…(5) Its solution
(iii) Series (4) and (5) are convergent at every
point within the circle of convergence at a.
[
y = a0 x – n 1 –
4(–
1
n + 1)
x2

Thus a solution in series will be valid only if

the series is convergent.
+ 2
1
4 .2! (– n + 1) (– n + 2) ]
x4 – …

General Method Bessel functions, Jn (x)

∞
(– 1)r
To solve in series the equation (1) when x = 0
is a regular singularity
Jn (x) = Σ
r=0 r ! (n + ()
r)!
x n + 2r
2
 Physical Sciences | 73

If n = 0 when n is odd last term

(– 1)r n.(n – 1) … 3.2
J0 (x) = Σ
(r !)2 ()x
2
2r
(–1) 1/2(n – 1)
(2n – 1) (2n – 3) … (n + 2)·2·4…(n – 1)
x
∞ So Pn (x) is Legendre’s function of the first
(– 1)r
& J– n(x) = Σ kind, Another solution
r=0 r! 
√ (– n + r + 1)
n!
Qn (x) =
· (x2) – n + 2r
1.3.5 … (2n + 1)

General solution of Bessel equation [x– n – 1 +

(n + 1) (n + 2) – n – 3
2.(2n + 3)
x +… ]
y = AJn (x) + BJ– n (x) The series for Qn (x) is a non-terminating
Recurrence relation : series.
(1) x J′n = n Jn – x Jn + 1 General Solution of Legendre’s Equa-
(2) x J′n = – n J n + x Jn – 1 tion
(3) 2J′ n = J n – 1 – Jn + 1 y = APn (x) + BQn (x)
(4) 2n J n = x (Jn – 1 + Jn + 1) Rodrigue’s Formule
d 1 dn 2
(5) (x – n. Jn ) = – x – n Jn + 1 P n (x) = (x – 1)n
dx 2n n! dx n
d Generating Function of Legendre
(6) (x n Jn) = x n Jn – 1
dx Polynomial
( ) 4 2
(7) J′ 2 (x) = 1 – 2 J 1(x) + J 0 (x)
x x
P n (x) is the coefficient of z n in the expansion
of (1 – 2xz + z 2 ) – 1/2 in ascending powers of z.
Where Jn (x) is the Bessel function of Ist kind.
Orthogonality
Orthogonality of Bessel Functions +1
1 ∫ Pm(x) Pn(x) dx = 0 n≠m
∫ x Jn (αx) . Jn (βx) dx = 0 –1
0 2
= n=m
where α and β are the roots of Jn (x) = 0 2n + 1
A Generating Function for Jn (x) Recurrence Relation
Jn (x) n
is the coeff. of z in the expansion of (1) nPn = (2n – 1)x Pn – 1(x) – (n – 1) Pn – 2 (x)
e x/2 (z – 1/z) (2) xP′n – P′n – 1 = nPn
Legendre’s Equation (3) P′n – xP′n – 1 = nPn – 1
(4) P′n + 1 – P′n – 1 = (2n + 1)Pn
d 2y dy
(1 – x 2 ) – 2x + n(n + 1)y = 0 (5) (x 2 – 1)P′n = n [x Pn – P n – 1]
dx 2 dx
Legendre Polynomial Pn (x) : (6) (x 2 – 1) P′n = n (Pn – 1 – x Pn)

P n (x) =
1.3.5. … (2n – 1) Laguerre’s Differential Equation
n! d 2y dy
x + (1 – x) + ny + 0
xn –[ n (n – 1) n – 2
(2n – 1).2
x +… ] dx 2 dx
Its solution given by Laguerre’s polynomial
dn
Legendre Function of IInd Kind Ln (x) = e x n [x n e – x]
dx
Actually solution for Pn (x) is a terminating
series. When n is even, its lest term being. Generating Function
∞
n (n – 1) (n – 2) … 1 e – xt/(1 – t) Ln (x) n
(– 1) 1/2n = Σ t
(2n – 1) (2n – 3) … (n + 1) 2.4.6 … n 1–t n=0 n!
74 | Physical Sciences

Orthogonal Properties where the Fourier coefficients an and b n can be

∞
taken as
∫ 0 e – x Lm (x).Ln (x) dx = {(n0!)2 mm ≠= nn
∫ – L f (x) cos nπx
1 +L
an = dx
L L
Recurrence Relation
(1) x L′′n (x) + (1 – x) L′ n (x) + nLn (x) = 0 and bn =
1 +L
L –L
∫
f (x) sin
nπx
L
dx
(2) (1 + 2n – x) Ln (x) + n 2 Ln – 1(x) – Ln + 1(x) If C is any real number the limits – L and + L
=0 can be replaced by C and C + 2L,
(3) L′n (x) = n L′n – 1(x) – n L n – 1(x) For L = π the series and the coefficients, an
and b n , are simple. f (x) is periodic for 2π. By
Hermite Polynomial putting x = 0 we can determine a0, thus a0 will be
Hermite’s diff. equation
∫ – L f (x) dx;
1 +L
y′′ – 2xy′ + 2ny = 0 a0 =
L
Rodrigue’s Formule a0
∫ – L f (x) dx
1 +L
=
dn 2 2L
Hn (x) = (– 1)n e x2 (e – x2)
dx n a0
is the mean of a f (x) over a period.
Generating Function 2

2
∞
Hn (x) n Odd and Even Functions
e 2tx – t = Σ t
n=0 n! A function f (x) is said to be odd if f (x)
= – f (–x) and even if f (x) = f (– x). For example
Orthogonal Property f (x) = cos x is an even function whereas f (x) = sin x
∞ ⎧ is an odd function. The figures given below
0 m≠n
∫ – ∞ e – x2 Hm (x) Hn (x) dx = ⎨⎩2n n! √⎯ π m=n
illustrate the example of odd, even and neither odd
nor even functions respectively.
Recurrence Relation
(1) H′n (x) = 2n Hn – 1 (x)
(2) 2x Hn (x) = 2n Hn – 1 (x) + Hn + 1 (x)
(3) H′n (x) = 2x Hn (x) – Hn + 1 (x)
(4) H′′n (x) – 2x H′n (x) + 2n Hn (x) = 0
Fourier Analysis f(x)
If any function satisfies the condition
f (x + T) = f (x) for each value of x, the function is f (x) = f (–x)
said to be periodic for time period T. (T > 0). e.g., –X +X
f (x) = cos x is a periodic function of time period
2π, 4π, 6π…etc. for it T = 2π = least period.
(even function)
Fourier Series
Let any function f (x) is defined in the interval
ranging between – L to + L (the period is 2L), and
outside this interval, f (x + 2L) = f (x). Then this f (x) ≠ ± f (–x)
function can be expand in terms of Fourier series.
∞
a
f (x) = 0 + Σ
2 n=0 ( an cos
nπx
L
+ bn sin
nπx
L ) (neither odd nor even function)
 Physical Sciences | 75

Half Range Fourier ‘sine’ or ‘cosine’ z = x + iy, then –z , known as complex

Series –
conjugate of z, will be z = x – iy such that
A series in which only sine or cosine terms
are present (respectively) the half-range series of a r =  z–z
√
function can be obtained by defining the function In polar form,
only in half-interval (and hence half range series), z1 z2 = (r1 eiθ 1 ) (r2 eiθ 2 )
(0 to L), which is half of (– L, L), and stating if
= r1 r2 ei(θ1 + θ2 )
the function is odd or even in the interval (0, L).
The function should be clearly defined in the other z1 r
Likewise = 1 ei (θ1 – θ2 )
half also. In such cases, z2 r2
2 L If n is any real number,
∫
nπx
an = 0 and bn = f (x) sin dx zn = (x + iy)n
L 0 L
for half range sine-series, = rn (cos θ + i sin θ)n
2 L = rn einθ
∫
nπx
bn = 0 and an = f (x) cos dx Also
L 0 L
for a half range cosine-series. z1/n = (x + iy)1/n
= r1/n e(1/n) iθ
Complex Analysis
θ + 2πk θ + 2πk
A number which consist of a real and an
imaginary part is said to be complex number. In
= r1/n cos [ n
+ i sin
n] ,

terms of two real numbers x and y, the complex k = 0, 1, 2,…n – 1

This is known as Morvie theorem and e± iθ
form is given as x + iy where i = √
– 1. (Remember, = cos θ ± i sin θ are known as Euler’s formulae.
y is imaginary part not iy). Hence two complex
numbers are said to equal if and only if there real Now, we have
and imaginary parts are equal separately. w = f (z)
The complex number can also be represented = f (x + iy)
in polar co-ordinates (r, θ) instead of Cartesian co- ⇒ w = u (x, y) + iν (x, y)
ordinates (x, y)
x = r cos θ w = u + iν like z = x + iy.
y = r sin θ w can be separated into real and imaginary
parts just as z can be separated into real and
or z = x + iy imaginary parts.
= r (cos θ + i sin θ) For example,
= r eiθ w = f (z) = 5z2
r is the modulus of absolute value, i.e., = 5 (x + iy)2
|r|=√
x2 + y2 . The graphical representation of a = 5x2 – 5y2 + 10ixy
complex number (in the complex plane) is shown = u (x, y) + iv (x, y)
below. where u = 5(x2 – y2) and v = 10xy.
y
z = reiθ Analytic Function
We assume that w is single valued, we have
r
dw df(z)
iy = = f ′(z)
dz dz
θ
x If f(z) exists at z = z0, we say that f(z) is
O x
analytic at z = z0 . If f(z) exists at everywhere in the
We can have a function of complex variable complex plane inside a region R, we say that f(z)
such that w = f (z). θ is called argument of z i.e., is analytic in the region R. For analyticity f(z)
θ = arg z. must be continuous single valued.
76 | Physical Sciences

Analytic functions sometimes also known as

holomorphic or monogenic functions. ∫
O f(z) dz = 0
C
The following two theorems can tell us This is called contour integral in complex
whether a function is analytic or not ? analysis.
Theorem 1. If f(z) = u (x, y) + iv (x, y) is Morera’s Theorem
analytic in a region then
∂u ∂v Morera’s theorem is another form of
= Cauchy’s integral theorem this states that if a
∂x ∂y
∂u ∂v function f(z) is continuous in a simply connected

and =–
∂y ∂x
The above set of conditions known as
Cauchy-Riemann conditions are necessary and
sufficient conditions for a function to be analytic.
Theorem 2. If f(z) = u + iv is analytic in a
region, then u and v satisfies Laplace’s equation in
the region. In this case u and v are called harmonic
function.
Again by differentiating above equations, we
∫
region R and O f (z) dz = 0 for every closed
C
contour C within R, then f(z) is analytic through-
out R.
Cauchy’s Integral Formula
If f(z) is a function, which is analytic on and
inside c closed contour C and z0 be any point
inside the contour then the value of f(z) at z0 i.e.,
f(z0 ) is given by

∫
get 1 f(z)dz
∂2 u ∂2 u f(z0 ) = O
= ; 2πi C (z – z0 )
∂x∂y ∂y∂x
the contour C is traversed clockwise always.
∂2 v ∂2 v Cauchy’s formula can be modified for nth deriva-
=
∂x∂y ∂y∂x tive too.
∂2 u ∂2 u

∫
Hence + = 0 n f(z)dz
∂x2 ∂y2 fn (z0 ) = O
∂2 v ∂2 v 2πi C (z – z0 )n + 1
and + = 0
∂x2 ∂y2 If a function of complex variable has first
Hence u and v satisfies Laplace’s equation for derivative it must have all higher derivatives.
two dimensions. (This is not a necessary condition for real deriva-
tives.)
Complex Integrals or Contour Integrals
Taylor’s Expansion
A continuous Jordan curve consisting of chain
of finite numbers of regular arcs is known as f(z) is to be analytic on and inside of a closed
contour. Thus a contour is smooth curve having circle C having its centre at z = a, then for all
no multiple points. points, z, in the circle we have
∞
Let f(z) is a continuous and single valued f(z) = Σ an (z – a)n
function in the region R, we can take the integral n=0

of f(z) along a path C from a given point z1 , to z2 where z is a point in C such that |z – a| < r and
where z1 = x1 + iy1 and z2 = x 2 + iy2 fn (a)
an =
and dz = dx + idy n
∫C dz = ∫C (x + iy) dx+ idy
f(z) fn (a) being the nth derivative of f(z) at z = a.
∞
(z – a) n n
w = f(z) = u + iv hence Hence f(z) = f(a) + Σ f (a)
∫
x2 y2 n=1 n
∫Cf(z) dz = x1 y1
[(udx – vdy) + i (vdx + udx)]
which is known as Taylor’s expansion
Cauchy’s Theorem f(z) = f(a) + f′(a) (z – a)
(z – a) 2 (z – a) 3
If f(z) is analytic on and inside a closed + f′′(a) + f′′′(a) +…
contour C, then Cauchy’s theorem states that 2 3
 Physical Sciences | 77

Laurent Theorem = (z – a) m [am + am + 1 (z – a)

If f(z) be analytic in the ring shaped region D + …]
bounded by two concentric circles C1 and C2 with ∞

centre z0 and radii r1 and r2 (r1 > r2 ) and z be any = (z – a) n Σ am + n (z – a)n

n=0
point in D then
∞ ∞ f (z) = (z – a) m φ (z).
f(z) = Σ an (z – z0 )n + Σ bn (z – z0 )– n with φ(z) = Σam + n (z – a)n
n=0 n=0

∫
1 f(z)dz = am + am + 1 (z – a) + …
where an = O is analytic and non zero in the neighbourhood of z
2πi (z – z0 )n + 1
= 0.
∫
1 f(z)dz
and bn = O Note—The zero is said to be small if m = 1
2πi (z – z0 )– n + 1
i.e. z = a is a simple zero of the function if f (a) =
n = 1, 2, 3,…… 0 and f ′(a) ≠ 0
If all b’s are zero, f(z) is analytic at z = z0 and
we call z0 a regular point. Singularities of an Analytic Function
If b n ≠ 0 but all b’s after bn are zero, f(z) is A point at which a function f (z) is not analytic
said to have a pole of order n at z = z0 if n = 1, we is known as singular point or singularity of f ( z).
say that f(z) has a simple pole. There are different types of singularity.
If there are infinite number of b’s different (1) Isolated and Non-isolated Singularity—
from zero, f(z) has an essential singularity at z = z0 . If z = a is a singularity of f (z) and if there is no
other singularity within a small circle surroun-
1 ding the point z = a, then z = a is said to be an
The coefficients b 1 of is called the
(z – z0 ) isolated singularity of the function f (z); otherwise
residue of f(z) at z = z0 . it is called non-isolated.
Zeros of an Analytic Function (2) Removable Singularity—Let f ( z ) be
analytic in Domain D except at the point z = a,
A zero of an analytic function f (z) is a value then z = a is an isolated Singularity of f (z) and we
of z such that f (z) = 0. can draw two concentric circles with centre ‘a’
If the function is analytic at a point a, then such that f (z) can be expanded in a Laurent’s
there exists a neighbourhood of the point z = a, series and we shall write
throughout which the function represented by a ∞ ∞
bn
Taylor series f (z) = Σ an (z – a)n + Σ n – (I)
n=0 n = 1 (z – a)
∞
∞
f (z) = Σ an (z – a)n In (I) Σ an (z – a)n is called the regular part
n=0
n=0
f n (a) ∞
bn
where, an =
n!
where z = a is a zero of f (z) then of the f (z) and Σ is called its principal
n=1 (z – a) n
f (a) = 0 part in the neighbourhood of a. If an = 0 for all n,
a0 = f (a) = 0 then the singularity at ‘a’ is called a removable
In an addition to this f ′(a) = 0 = f′′(a) = … = singularity of f (z).
fm – 1(a), but fm (a) ≠ 0, then f (z) is said to have a (3) Essential Singularity—In equation (I)
zero of order m at z = a. It may also seen that since above, if the principal part contains infinite
f m (a) , number of terms then z = a is called an isolated
am = hecne at a zero of order m, we have a0
m! essential singularity of f (z).
= 0, a1 = a 2 … = a m – 1, but a m ≠ 0. Thus in the ∞
neighbourhood of the zero z = a of order m the (4) Poles—In series (I) above, Σ an (z – a) n
n=1
function f (z) is represented as the series.
consists of finite number of terms. Thus z = a is
f (z) = am (z – a)m + am + 1 (z – a)m + 1 said to be a pole of order m of the function f (z).
+… when n = 1, the pole is said to be simple.
78 | Physical Sciences

Residues n
The coefficients, in the Laurent’s series for (2) L(t n) = , when n = 0, 1, 2, 3.
sn + 1
f(z), can be obtained in a customary manner, by
writing the coefficients for the Taylor’s series, 1
(3) L(e at) =
corresponding to (z – a ) n f (z). The coefficients s–a
a–1 is called residue of f(z) at the pole, z = a. This a
(4) L(sin at) = 2
is of prime importance. s + a2
s
1 dn – 1 n (5) L(cos at) = 2
a–1 = Lim z → a n – 1 [(z – a) f(z)] s + a2
n – 1 dz a
(6) L (sinh at) = 2
nth is the order of pole for a simple pole, i.e., s – a2
n = 1, s
(7) L (cosh at) = 2
s – a2
a–1 = Lim z → a (z – a) f(z)
∴ Properties of Laplace Transforms
( 0 =1 )
(I) Linearity property—If a , b , c be any
The Residue Theorem constants and f, g, h any functions of t, then
If f(z) is analytic, within and on the boundary L[af (t) + bg (t) – ch (t)] = aL{f (t)}
of a region R except at a finite number of poles, + b L {g (t)} – cL {h (t)]
a–1, b–1, c–1……etc. within the region and a–1, b– 1,
(II) First shifting property—
c–1…… etc. are the residues at these points then,

∫
—
O f(z)dz = 2πi (a–1 + b– 1 + c –1 + ...) If L{f (t)} = f (s), then
C —
L{e at f (t)} = f (s – a)
∫
i.e., the integral of f(z) i.e., O f(z)dz will be
C
Application of this property leads to
1
(2πi) times the sum of residues of f(z) at the poles (1) L(e at) =
enclosed by c. s–a
n!
C (2) L(e at t n) =
(s – a) n + 1
a a–1 b
(3) L(e at sin bt) =
(s – a) 2 + b 2
b s–a
b–1 (4) L(e at cos bt) =
(s – a) 2 + b 2
c c
–1 b
(5) L(e at sinh bt) =
(s – a) 2 – b 2
The integration is taken counter clockwise. s–a
(6) L(e at cosh bt) =
Laplace Transform (s – a) 2 – b 2
Denoted by L{f (t)} is defined by Where in each case s > a

∫
∞ (III) Change of scale property—
L{f (t)} = e – st f (t)dt —
0
If L{f (t )} = f (s),
Provided that the integral exists. S is a 1— s
parameter which may be real or complex number. then L{f (at)} =
a
f () a
Transforms of Elementary Functions Transforms of Derivatives
The direct application of the definition gives
the following formulae— (1) If f ′ (t) be continuous and L{f(t)} = f(s),
then
1 —
(1) L(1) =
s L{f ′ (t)} = s f (s) – f (0)
 Physical Sciences | 79

(2) If f ′ (t) and its first (n – 1) derivatives be Fourier Transfroms

continuous, then ∞
—
L{f n (t)} = s n f (s) – s n – 1 f (0)
If F(s) = ∫ – ∞ f (t) e ist dt …(1)
– s n – 2 f ′ (0) … – f n – 1 (0) ∞
∫ – ∞ F(s) e – isx dx
1
Then f (x) = …(2)
If L{f (t)} = f (s), then 2π
dn — The function F(s), defined by (1), is called the
L{t n f (t)} = (– 1)n n [ f (s)],
ds Fourier transform of f(x). The function f (x), as
where n = 1, 2, 3, …… given by equation (2), is called the inverse
fourier transform of F(s).
Inverse Transforms
Standard forms : Fourier Sine and Cosine Transform
(1) L– 1 []
1
s
=1 Fs (s) =
∞
∫ 0 f (x) sin sx dx
(2) L– 1 [ ]
s–a
1
= e at then f (x) =
2
π
∞
∫ 0 Fs (s) sin sx dx
tn – 1 ,
(3) L– 1 []
1
sn
=
(n – 1) !
n = 1, 2, 3, … and F c (s) =
∞
∫ 0 f (x) cos sx dx
e at t n – 1
(4) L– 1 [ ] 1
(s – a) n =
(n – 1) ! then f (x) =
2
π
∞
∫ 0 Fc (s) cos sx dx
(5) L– 1 ( ) 1
s2 + a2
1
= sin at
a Finite Fourier
Transforms
Sine and Cosine

(6) L– 1 ( ) s
s2 + a2
= cos at The finite fourier sine transform of f (x), in
0 < x < c is defined as
(7) L– 1 ( ) 1
s 2 – a2
1
= sinh at
a
∫
c
F s (n) = 0 f (x) sin
n πx
dx
c
(8) L– 1 ( ) s
s2 – a2
= cosh at
where n is an integer
2 ∞ n πx
(9) L– 1 [ ] 1
(s – a) 2 + b 2
1
= e at sin bt
b and f (x) = Σ F (n) sin
c n=1 s c
(10) L– 1 [ ]s – a
(s – a) 2 + b 2
= e at cos bt is known as inverse finite Fourier sine trans-
form of Fs (x) the finite Fourier cosine transform
(11) L– 1 [ ] s
(s 2 + a 2 )2
=
1
2a
t sin at of f (x), in 0 < x < c, is defined as
c
∫
nπx
Fc (n) = 0 f (x) cos dx,
(12) L– 1 [ ] 1
(s 2 + a2 )2
= 3
2a
1 c
where n is an integer whereas inverse finite
(sin at – at cos at) Fourier cosine transform of Fc (n) is
Laplace Transform of Special Functions 1 2 ∞ n πx
1 f (x) = Fc (0) + Σ Fc (n) cos
(1) L [J0 (x)] = c c n=1 c

√ s2 + 1
(2) L[J1 (x)] = – L [J′0 (x)] Properties of Fourier Transforms
= – [sL {J0(x)} – 1] (1) Linear property—If F(s) and G(s) are
s Fourier transforms of f (x) and g(x) respectively,
= 1–

√ s2 + 1
then
s F[af (x) + bg (x)] = aF(s) + bG(s)
(3) L{t J0 (at)} = 2
(s + a 2 )3/2 where a and b are constants.
80 | Physical Sciences

(2) Change of scale property—If F(s) is the Relation between Fourier and Laplace
complex Fourier transform of f (x), then Transformation
F[f (ax)] = F
1
a ()s
a
;a≠0 If f (t) = {e 0
– xt
g(t) ‚ t > 0
‚ t<0
(3) Shifting property—If F(s) is the complex Then F{f (t)} = L{g (t)}
Fourier transform of f (x), then
F[f (x – a)] = e isa F(s) Elementary idea of Tensor
(4) Modulation theorem—If F(s) is the Quantity associated with two or more
complex Fourier transform of f (x), then directions is known as tensor.
1 Summation Convention
F{f(x) cos ax} = [F (s + a) + F (s – a)]
2 Consider a sum of the type
Note—If Fs (s) and Fc (s) are Fourier sine and n
cosine transforms of f (x) respectively, then a1 x 1 + a2 x 2 + … + a nx n i.e. Σ ai x i
i=1
1
(i) Fs [f (x) cos ax] = [Fs (s + a) + Fs (s – a)] In tensor analysis, the subscripts of the
2
1 symbols x 1 , x 2 , … x n are replaced by superscripts
(ii) Fc [f (x) sin ax] = [Fs (s + a) – Fs (s – a)] and we write these as x 1 , x 2 , … x n . The
2
superscripts do not stand for the various powers of
1 x but act as labels to distinguish different symbols.
(iii) Fs [f (x) sin ax] = [Fc (s – a) – Fc (s + a)]
2
The power of a symbol (say : x i) will be
Convolution n
indicated as (x i)2, (x i)3 etc. … Σ ai x i
The convolution of two functions f (x) and i=1
g (x) over the interval (– ∞, ∞) is defined as The repeated index i over the summation is to
∞
f×g = ∫ – ∞ f (u) g (x – u) du = h (x) be done is called dummy index.
e.g. s = a ij x i x j takin n = 3
Convolution theorem for Fourier transforms :
first term sum on i from 1 to 3
The Fourier transform of the convolution of
f ( x ) and g (x) is the product of their Fourier s = a 1j x′xj + a 2j x 2j + a3j x 3j
transforms i.e., Now repeat index j from 1 to 3
F{f (x).g (x)} = F{f (x)}.F{g (x)} s = a11x 1x + a12 x 1x 2 + a13 x 1x 3 + a21 x 2x 1
Parseval’s Identity for Fourier Trans- + a 22 x 2x 2 + a23 x 2x 3 + a31 x3x 1
forms + a32 x 3x 2 + a33 x 3x 3
If the Fourier transforms of f (x) and g (x) are e.g. If f is a function of n variables x i write
F(s) and G(s) respectively, then differential of t
1 ∞ — ∞ —
(i)
2π – ∞∫ ∫
F(s) G (s) ds = – ∞ f(x) g (x) dx f = f (x 1, x 2, … x n)
∴ From calculus, we have
1 ∞ ∞
(ii) ∫
2π – ∞
|F(s) | 2 ds =
∫–∞
| f (x) |2 dx
df =
∂f ∂f
dx 1 + 2 dx 2 + … + n dx n
∂f
∂x 1 ∂x ∂x
Note—
∂f
2 ∞ ∞ dx i
∫ ∫
=
(i) F (s) G (s) ds = f (x) g (x) dx ∂x i
π 0 c c 0
Transformation of Co-ordinates
2 ∞ ∞
(ii) ∫
π 0
F s (s) G s (s) ds = 0∫ f (x) g (x) dx In 3D space → Co-ord of points are (x 1 , x 2 ,
x 3) referred to a particular frame of reference
2 ∞ ∞
(iii)
π∫ 0 ∫
[Fc (s)]2 ds = 0 {f (x)}2 dx n dim. space → (x 1 , x 2, … x n)
– – –
2 ∞ ∞ Let (x 1 , x 2 , … x n ) be the co-ord of the same
(iv)
π∫ 0 ∫
| F s (s) |2 ds = 0 |f (x) |2 dx
point referred to another frame of reference.
 Physical Sciences | 81

– – – Compactly
Suppose x1 , x 2 , … x n are independent single
valued functions of x , x 2, … x n so that
1 a 1j A1j = ∆,
–x 1 = φ1 (x 1, x 2, ……, x n) a 1j A2j = 0,
a 1j A3j = 0
–2
x = φ2 (x 1, x 2, ……, x n) a 1j A1j = ∆ δ1 k
··· a 2j Akj = ∆ δ2 k,
–n a 3j A kj = ∆ δ3 k
x = φn (x 1, x 2, ……, x n)
–i All nine components included in aij Akj = ∆δik
x = φ i (x 1, x 2, ……, x n) –
If x i & x i are independent co-ordinate of a
Similarly point show that
– – –
x i = ψ i (x 1 , x 2 , ……, x n ) –
∂x j ∂ x k j
Scalars or invariants—A function φ(x 1 , x 2 , – k ∂ x i = δi
x3) is called a scalar or an invarient of its original ∂x
⊇

value does not change upon transformation of Co- Contravarient Vectors

– – –
ord. from x 1 , x 2, x 3 to x 1 , x 2 , x 3 .
Let A1 , A2 , …… A n (i.e., A i) be a set of
– – –
φ(x 1, x 2, x 3) = ψ (x 1 , x 2 , x 3 ) n functions of Co-ord system x1, x 2 , … xn (i.e. x i).
A scalar or invariant is also called a tensor of If these transform in another system of Co-ord
order (or rank) zero. –1 –2 – –
x , x , ……, x n (i.e., x i) acc. to … law.
Kronecker delta—The quantity δij defined –
by the relation ∂ xi j
Ai = A
∂xj
δi j = 0 when j ≠ i
A → is component of a contravarient vector
i
δi j = 1 when j = i is called or contravarient tensor of order one
Kronecker delta. –
Evidently – ∂ x i 1 ∂ –x i 2
e.g., d xi = dx + 2 dx
δ1 1 = δ2 2 = δ3 3 = …… = δn n = 1 ∂x 1 ∂x
While δ1 2 = δ2 3 = δ3 2 = …… = 0 –
∂ xi
+ …… + n dx n
We note that by summing up with respect to, ∂x
the repeated index j –i
∂x
a3j δ2j = a31 δ21 + a32 δ22 + a33 δ23 =
∂x j
dx j
+ a34 δ24
= 0 + a32 + 0 + 0 = a32 Covariant Vectors
In general, Let A i be a set of n functions of the co-ord
a ij δk j = a i1 δk1 + a12 δk2 + … + a ik δ kk system x i. If these transform in another system of
—
+ … a in δk n co-ord x i acc. to law
= 0 + 0 + … + a ik.1 + … = a jk — ∂xj
A i = – Aj
Show that ∂ xi
a ij A kj = ∆ δ ik, Ai = are called covariant vector or tensor of
where ∆ is det of order 3 & A ij are Co-factors order as.
of a ij ∂φ ∂φ ∂x1 ∂φ ∂x 2
a11A11 + a12A12 + a13A13 = ∆ –i = ∂x 1 –i + ∂x 2 ∂x i
∂x ∂x
a11A21 + a12A22 + a13A23 = 0 ∂φ ∂x n
+ …… +
a11A31 + a12A32 + a13A33 = 0 ∂x n ∂x i
82 | Physical Sciences

e.g. A covariant tensor has components xy, Addition of Tensors

2y – z2 , xz in rect. Co-ord. Find its covariant
The sum (or difference) of two tensors of the
component in spherical Co-ord. same order and type is another tensor of the same
1 2
x = x, x = y, x = z 3
order and type
–1 – – ∂x k ∂x l
and x = r, x 2 = θ, x 3 = φ —
A ij = A
A1 = xy, A2 = 2y – z2 , A3 = x z — — kl
∂ xi∂ xj
Acc. to transformation — ∂x k ∂x l
— ∂x j and B ij = — — Bkl
Ai =
– j
A ∂ xi∂ xj
∂ xi — — ∂x k ∂x l
— — — A ij ± B ij = — — (Akl ± Bkl)
We want to evaluate A1, A2, A3 where A1, A2 , ∂ xi∂ xj
A3 are known — ∂x k ∂x l
x = r sin θ cos φ, C ij = C
— — kl
∂ x i∂ xj
y = r sin θ sin φ;
z = r cos θ Outer Product
— ∂x 1 ∂x 2
Now, A1 = – A 1 + – A 2 ij — —
∂x 1 ∂ x1 C kl = A ij B kl
∂x 3 ∂ —x i ∂ x j pq
—
=  p
+ A
–1 3 A 
∂x ∂x ∂x q
∂x ∂y ∂z  ∂x r ∂x s B 
= xy + (2y – z 2 ) + xz
∂r ∂r ∂r  —k —l rs
∂ x ∂ x 
= sin θ cos φ.(r sin θ cos φ).(r sin θ sin φ)
+ sin θ sin φ [2r sin θ sin φ – r2 cos2 θ] — —
∂ x i ∂ x j ∂xr ∂x s pq
= A Brs
+ cos θ.r sin θ cos φ.r cos φ ∂x p ∂x q ∂ — —
xk∂ xl
— — pq
Similarly A2 and A3. = …… C rs
Tensors of higher order Contraction of a Tensor
Let i & j be each given values 1 to n, then the — — —
symbol Aij will give rise to n 2 function
—ijk ∂ x i ∂ x j ∂ x k ∂x s pqr
Al = A
∂x p ∂x q ∂x r —i s
(1) If Aij be a set of n 2 function of Co- ∂ x
— put the covarient index l = a contravariant
ordinates x 1 , x 2, …, x n which transform to A ij in
index i, so that
another system of Co-ord
— —
—ij – –
∂ x i ∂ x j kl
—ijk ∂ x i ∂x j ∂ x k ∂x s pqr
Ai = Ap
A =
∂x k ∂x l
A ∂x ∂x ∂x r ∂ —
p q
xi
—ij — —
A = components of contravarient tensor of ∂ x i ∂ x k s pqr
= δ A
order 2. ∂x p ∂x r p s
— ∂ x k ∂x l — —
(2) If A ij = – – Akl ∂ x j ∂ x k pqr
∂ xi ∂ xj =
∂x q ∂x r p
A
→ Covariant tensor of second order Contravarient tensor of order 2.
–
—i ∂ x i ∂x l k
(3) If Aj = A Rotation of Axis and Vectors
∂ xk ∂ –x j l
— — — —
→ Mixed tensor of the second order A = i A1 + i 2 A2 + i 3 A3
 Physical Sciences | 83

Let the set of axis be rotated by keeping the = Ai1 A j 2 A k3 ∈123 + Ai1 Aj3 Ak2 ∈132
origin O fixed, leading to a new axis 1′, 2′, 3′. In + Ai2 Ai3 Ak1 ∈231 + Ai2 Aj1 Ak3 ∈213
— — —
terms of new vectors i′1 , i′2 , i′3 the same vector + Ai3 Aj1 Ak2 ∈312 + Ai3 Aj1 Ak2 ∈312
can be expressed as + Ai3 Aj2 Ak1 ∈321
— — — — = Ai1 Aj2 Ak3 ∈123 – Ai1 Aj3 Ak2 ∈132
A = i′1 A′1 + i′2 A′2 + i′3 A′3
+ Ai2 Aj3 Ak1 ∈231 + …
— — — — —
⇒ i′1 A′1 + i′2 A′2 + i′3 A′3 j ∂x l ∂ x j m ∂ x j ∂x m
(1) δi = δ =
— ∂x m l — —
= i 1 A1 + i 2A2 + i 3A3 ∂xi ∂ xm∂ xi
1
⇒ taking scalar product with i 1 on both side ∂xj
—
j m
= — = δi [δi = 0 for l ≠ m]
—— ——
1
A1 = (i′1 . i1 )A 1 + (i′1 i 2 )A 2 ∂ x i

— dx1 , dx 2 , dx 3
+ (i′1 .i 3)A 3 (2) be the component of fluid
dt dt dt
— dx i
i′1 .l 1 = a11, velocity of the point (x 1 , x 2, x 3) i.e. be the
dt
i′1.l 2 = a12, component of velocity in the co-ordinate system x i
— — —
—
i′1 .l 3 = a13 ∂xi ∂ x j dx 1 ∂ x j dx 2
= +
dt ∂x1 dt dx 2 dt
Similary —
A′ 1 = a11A1 + a12A2 + a13A3 ∂ x j dx 3
3
+ 3
∂x dt
A′ 2 = a21A1 + a22A2 + a23A3
—j —j i
A′ 3 = a31A1 + a32A2 + a33A3 ∂x ∂ x dx
= → tensor of rank(1)
Where aim = i′ i.i m dt ∂x i dt
3
—i
(3) x = a im x m + b i …(1)
A′ 1 = Σ aij Aj i = 1, 2, 3, …
∂x i
j=1
⇒ = a im …(2)
A set of nine quantities ∂x m
B ij, i = 1, 2, 3 & j = 1, 2, 3, form a tensor of Also from (1),
rank 2. —
ai r x i = a i r a i m x m + a i r b i
3 3
B′ij = Σ Σ ail ajm Blm = δr m x m + a ir b i
l=1m=1 = x r + a ir b i
Three order ∂x r
i
3 3 3
— = ar
C′ijk = Σ Σ Σ ail ajm akn Clmn ∂xi
l=1m=1 n=1
∂x m
Rank of order 3. i.e. = a im
∂x i
Kronecker delta δij
δ′ij = δij (isotropic tensor) Group Theory
The Levi Civita symbol Definition
∈ijk = 0 if any two of the indices are equal A group G is a collection of elements with a
= 1 if ijk is an even permutation of 1, 2, 3 rule to combine two elements given. This process
of combination is called group multiplication
= – 1 if ijk ” ” odd permutation of 1, 2, 3 which must satisfy the following properties :—
Aij → 3 × 3 matrix (1) Closure—If A and B are any two
∈ lmn → Zero is any two indices are equal elements in G, then AB is an element of G.
84 | Physical Sciences
(2) Associativity—A(BC) = (AB)C, A, B, C
being any three elements in G.
(3) Existence of the Indentity element—
There is an element E in G such that AE = A,
where A is any element in G, E is called the
identity element.
(4) Existence of Inverse—For every element
there exists another element (denoted by A– 1) in
G, Such that AA – 1 = E. A– 1 is called the inverse
of the element A.
EA = A
and A– 1 A = E
By definition EE = E
Now, multiplying the relation AA– 1 = E from
the left by E, we obtain
or
EA– 1A = EE
(EA)A– 1 = E.
Thus EA = A
Similarly, multiplying AA– 1 = E by A from
the right one obtains
AA–1A = EA
or A(A– 1A) = A
Thus A– 1 A = E
In general group multiplication is not
commutative. A commutative group is known as
an abelian group.
For a finite group, the number n of element is
called the order of the group.
We now consider some examples of groups.
Example 1. The Four number 1, i, – i, – 1,
where i = √ – 1 form a group under ordinary
multiplication. If we multiply any two elements,
the result is one of the four elements. Thus
closure is satisfied. The multiplication is
associative. The element is obviously the identity
element. Also.
i (– i) = 1 and (– 1) (– 1) = 1
Showing that – i is the inverse of the element
i and vice versa and – 1 is its own inverse.
If we denote the element i by A, the elements
of the group can be specified as A, A2, A3, A4 = E.
Such a group, in which all the elements are
powers of a single element is called a cyclic
group. This is an example of a cyclic group of
order four. A cyclic group of order n is denoted
by cn.
In example, group multiplication was
ordinary multiplication. This need not be the case.
The set of all positive and negative integers
including zero form a group under addition. The
addition of two such integers gives another
integer. Such addition is associative. The identity
element in this case is O, because N + O = N,
Also, for every integer N, – N which also belongs
to the set, is its inverse : N + (– N) = O
The same set of positive and negative
integers, however, do not form a group under
ordinary multiplication. The product of two such
integers is another integer. The multiplication is
associative. The integer 1 is the identity element.
But, for a given integer N, its inverse now would
be N– 1, which is not an integer, if N ≠ 1.
3. The four matrices
( )
E = 0 1 ,
1 0
( )
A = –1 0 ,
0 1
B = (
0 – 1)
–1 0
,
C = ( 1 0)
0 –1
Form a group under matrix multiplication. It
can be verified by direct multiplication that the
product of any two of them gives one of the four
matrices. Matrix multiplication is associative.
( ) 1 0
The unit matrix 0 1 is obviously the
identity element.
Also, AC = CA = E,
B 2 = E,
indicating that A and C are the inverses of each
other and B is its own inverse.
4. The set of all n × n unitary matrices forms a
group under matrix multiplication. If U 1 and U 2
are two n × n unitary matrices, then U = U1U2 is
also a n × n unitary matrix.
U+U = (U 1 U2)+ U1 U2
= U2+ (U1 + U1 ) U2
= U2 + U2 = I
The multiplication is associative. The n × n
unit matrix, which is also a unitary matrix, is the
identity element. Also, for a unitary matrix U its
adjoint U+ is the inverse. But U+ is also a unitary
matrix :
(U +)+ U+ = UU+ = 1
and U+(U +)+ = U+U = 1
 Physical Sciences | 85

5. The set all symmetry transformations Thus to find the effect of AC on the
which leave a physical system invariant forms a equilateral-triangle, let us label the corners a, b, c.
group. Then
Group multiplication in this case is the a a b
application of two transformations successively. If
each transformation leaves the system invariant, AC = A =
the product will also leave the system invariant.
The product is thus a symmetry transformation.
The successive application of three transfor- b c c b a c
mations obeys associativity. The identity element Since C is reflection at axis 1, the effect of C
is the transformation which leaves the system on the triangle is to interchange b anc c. The
undisturbed. operation A then rotates the triangle about the z-
Finally, for every symmetry transformation axis to give the final configuration shown above.
there is the inverse operation. (for example, the But this is equivalent to Thus AC = ≠
π a
inverse of the operation “rotation by in the
2
anticlockwise direction about the z-axis’ will be F
π
“rotation by about the same axis in the
2
clockwise direction”). The inverse operation will
also leave the physical system invariant and is b c
thus a symmetry operation. The rules of multiplication for a finite group
As a specific example, consider the group can be expressed by multiplication table.
of symmetry transformations of an equilateral Multiplication table for the symmetry group of an
triangle. (This group is known as dihedral group, equilateral triangle.
D3). E A B C D F
We list the symmetry operation with the E E A B C D F
corresponding symbols below : A A B E F C D
Y B B E A D F C
C C D F E A B
D D F C B E A
3 2 F F C D A B E
6. The permutation of n identical objects form
X a group. This group is denoted by Sn . Each
permutation can be considered a transformation on
the system. Two successive such operations is
1
equivalent to another permutation. As a specific
case let us consider S 3 , the group of permutations
E : No transformation of three identical objects. A convenient notation to
2π denote a permutation is as follows :
A : anticlockwise rotation by about the z-
3
axis
4π
( 1 2 3
)
2 3 1 denotes a permutation in which
B : Anticlockwise direction about the z-axis object 2 has been brought to the first position,
3 object 3 to the second position and object 1 to the
C : reflection at axis 1 third position. The first row, therefore, indicates
D : reflection at axis 2 the positions and second row refers to the objects.
F : reflection at axis 3 For three objects, 6 permutations are possible. Let
It is very easy to verify that a combination of us lable them, as
any two of the operations is equivalent to one of
the above operations.
1 2 3
E = 1 2 3 , ( )
86 | Physical Sciences

between the elements : Gi ↔ G′i such that Gi Gj =

A = (13 2 3
1 2 ) ; Gk implies G′i G′j = G′k and Vice Versa.
If two groups are isomorphic to each other
B = ( )
1 2 3
; they have the same number of element and the
2 3 1
same multiplication table.
C = (1 )
1 2 3
3 2 ; Example of this is the dihedral group D3 and
the permutation group S3. The group (1, i, – 1, – i)
D = (3 )
1 2 3 is isomorphic to the cyclic group (A, A2 , A3 ,
2 1 ;
A4 = E)
F = (2 )
1 2 3 A homomorphic mapping of a group G to
1 3 another group G′ means a mapping (not
Then one can work out the effect of two necessarily one to one) such that G i Gj = Gk
successive permutations implies G′ i G′j = G′k.

( 1 2 3
AF = 3 1 2 )( 1 2 3
2 1 3 ) In homomorphism, in general, several
elements of G are mapped into one element of G′,
such that multiplication is preserved.
( 1 2 3
= 3 2 1 =D ) Isomorphism is a special case of homomor-
phism.
In this way, if the multiplication table of the
group is worked out, it will be found that the table Group Representation
is same as for the dihedral group, D3 ! If there is a homomorphic mapping from a
^
group G = {Gi} to a set of operators D(G) on a
Rearrangement Lemma
In a group multiplication table, each element vector space V, then the set of operators D(Gi) is
occurs only once in a row or column. said to form a representation of a group in the
D3 table is true in general : vector space V.
By a representation of a group G = {Gi} is
An element can occur twice in a row (or
column) if it is the result of two different multipli- meant a homomorphic mapping of G to a set of
cations. In a given group, let AX = Y & AB = Y. n × n non-singular matrices, D(Gi). This means
Then AX = AB. Multiplying from the left by A– 1, that for every element G i, we have an n × n non
we get AX = B. This is a contradiction because all singular matrix D(Gi).
the elements are supposed to be distinct. Gi G j = Gk
Sub group—For every group there are two If D(Gi) D(G j) = D(Gk)
trivial subgroups :— The identity element corresponds to the n × n
(i) the identity element E by itself forms a unit matrix.
subgroup. If Gk = E, the identity element Gi and Gj are
(ii) also the whole group G can be considered inverse of each other. Then D(Gj) D(Gj) = 1 & the
as a subgroup of itself. matrix D(Gi) and D(G j) are inverse each other.
A subgroup H of G is called a proper sub- The dim of matrices is called dim of
group if H ≠ G. representation.
Def—If A and X are any two elements in G Example 1. Consider the cyclic group C4 =
then XAX– 1 belongs to G. The elements A and (A, A2 , A3, A4 = E) and the set of matrices
XAX – 1 are said to be conjugate to each other.
A subgroup H is called invariant sub group, ( ) 0 1
D(A) = – 1 0 ,
if all its conjugate sub groups in G are identical
with H. D(A2 ) = ( )–1 0
0 –1 ,
Homomorphism and Isomorphism
Two groups G = (G1 , G2 … Gn ) and G′ =
( )
D(A3 ) = 1 0
0 –1

(G′1, G′ 2 , … G′n) are said to be isomorphic to each

other, if there is a one to one correspondence
and ( )
D(E) = 0 1
1 0
 Physical Sciences | 87

These matrices form a group under multipli- D(B) = 1,

cation & there is an isomorphism between the D(C) = – 1,
group and C4,
D(D) = – 1,
This set of matrices form a 2 dim faithful (If D(F) = – 1
mapping is isomorphic then representation is
called a faithful representation) representation of The following set of 2 × 2 matrices from a
group C4. two dim faithful representation of the group.
But this is not the only representation of C4 .
We can also have the following one-dim
D(E) = ( )
1 0
0 1
,

representation – 1 √ 3
D(E) = 1,  2 –

2
D(A) =
 √ 3 
 2 
D(A) = – 1, 1
–
D(A2 ) = 1, 2

– 1 √ 3
D(A3 ) = – 1
Another one-dim representation
2 2 
D(E) = 1, D(B) =
– √ 3 –1 
D(A) = 1,  2 2 
D(A2 ) = 1,
D(A3 ) = 1 D(C) = (– 01 – 10)
 1 √ 3
Since several representations of a group are
possible, representations are distinguished by
giving a superscript; D(µ) (A) etc. 2 2
D(D) =
√ 3 – 1
Thus the two one-dim representation given
can be denoted by D(1) & (D2) respectively. 2 2
We have two additional one dim represen-
 1 – √ 3
 2 2
tation of C4 :
D(3) (E) = 1, D(F) =
– √ 3 – 1 
D(3) (A) = i,  2 2 
D(3) (A2 ) = 1,
D(3) (A3 ) = – i Reducible and Irreducible Representa-
& D(4) (E) = 1, tion
D4 (A) = – i, Two representations D(Gi) and D′(Gi) are said
D (A2 ) = – 1,
(4) to be equivalent if D′(Gi) = S– 1 D(Gi) S for all G i
D(4) (A3) = i where S is a non-singular matrix.
So the dim of the representation has nothing A representation is said to be reducible, if it is
to do with the order of the group. Thus, C4 is a equivalent to a representation
D(N)(G i) A(Gi)
group of order four but it one or two dim
representation. Higher dim representation of C4 D′(Gi) = ( 0 D(V)(G i) )
are also possible. In other words a representation is reducible if
Example 2. For a dihedral group, D3 = (E, A, all the matrices in the representation can be
B, C, D, F) we have one dim representation D(E) brought to the above block diagonal form, by the
= 1, D(A) = 1, D(B) = 1, D(C) = 1, D(D) = 1, D(F) same similarity transformation.
= 1. If no similarity transformation can bring all
& D(E) = 1, the matrices to the block diagonal form, than the
D(A) = 1, representation is said to be irreducible.
88 | Physical Sciences

A representation is said to be fully reducible

if the representing matrices can be brought to from D′(A3) = (D 0(A ) D 0(A )),
(3) 3
(4) 3
the block diagonal form
D′(E) = ( 0 D (E))
D (E) (3)
0
D(N) (G i)
D(Gj) = ( 0
0
b (V) (Gi) ) (4)

Thus, the given two dim representation of C4

where, the representation D(N) (Gi) and D (V) is reducible and can be decomposed into the one
(G i) are irreducible. One can then say that dim representation D ( 3 ) and D ( 4 ) . This is
representation D(Gi) decomposes into irreducible symbolically stated as D′ = D (3) + D (4).
representations D(µ) (G j) and D(V) (G i)
Example 3. We have considered the The Regular Representation
following four dim representation of C4 : The group multiplication table is rearranged
such that the identity element E occurs along the
D(A) = ( ) 0 1
–1 0
, diagonal.
By Definition the matrix D(R) is suh that
D(A ) = ( 0 – 1),
2 –1 0
Dij (Gk) = 0 if G– 1i Gj ≠ Gk
= 1 if G– 1i Gj = Gk
D(A ) = (1 0),
3 0 –1
G1 = E G2 G3 ……
D(E) = (0 1)
1 0 G1 = E G2 G3 ……
–1
Let us transform all these matrices by the G1
similarity transformation S– 1 D(Gi) S, –1
G2

Where S =
1
( 1 1
 2 i –i
√
) –1
G3

D′(A′) = S – 1 D(A) S The Orthogonality Theorem

If D(µ ) and D(v ) are two irreducible
=
2
√
1
( 1 i) (– 1 0),
1 –i 0 1
representation of group G of order g, then
G
(µ) (v) –1 g
1
2
√
( 1
i –
1
i ) Σ Dil (Gk) Dmj (Gk ) =
k=1
δ δ δ
nµ µv ij lm
where n µ is the dimension of the
= (0 – i )
i 0 representation D (µ).

Similarly, Elements of Compulational Techniques

(NA)
D′(A2) = (– 01 – 10); Numerical Solution of Algebraic
D′(A ) = ( 0 i );
3 –i 0 Equation or Roots of Equations
Given an equation f (x) = 0
D′(E) = (0 1)
1 0 (1) Iteration method—When there is no
and
formula for exact solution, we use this method.
Obviously (2) Newton-Rapson method f (x) = 0—
D′(A) = (D 0(A) D 0(A)),
(3)
(4)
f (x) = 0, where f is assumed to have a continuous
derivative f ′
f (x n )
D′(A ) = ( 0
D (A )),
D (A ) (3)
0 2
2 xn + 1 = xn –
(4) 2 f ′ (x n)
 Physical Sciences | 89

(3) Secant method (Regular falsi)— x! , (n < x)

Replacing the derivative f ′ (x n ) by the difference x (n) =
(x – n) !
quotient, in Newton’s method. ∆x n = nhx (n – 1)
∆2 x (n) = n (n – 1)h 2 x (n – 2)
y = f (x)
··
∆ n x n = n ! h 2n
(E) Backward differences—
∇ f (x) = f (x) – f (x – h)
xn + 1 xn − 1 ∇2 f (x) = f (x – 2h) – 2f (x – h) + f (x)
··
f (x n) – f (x n – 1) n
So f ′ (x n ) = ∇n f (x) = Σ (– 1) n – r n C r f (x – rh)
xn – xn – 1 r=0
xn – xn – 1 (F) Relation between ∇, ∆ and E—
x n + 1 = x n – f (x n ) ·
f (x n ) – f (x n – 1) ∇E = E∆ ≡ ∆
(4) The bisection method—Suppose f is a E– 1 = 1 – ∇
continuous function defined on the interval [a, b], E = (1 – ∇)– 1
with f (a ) and f (b) of opposite sign. By the
intermediate value theorem, there exists a number (G) Central difference—
p in (a, b), with f (p) = 0. (a) Central difference operator δ—
[Important—This procedure will work for
the case where f (a) and f (b) will have opposite ( ) ( )
1
(i) δ f (x) = f x + h – f x – h
2
1
2
signs and there is more than one root in the (ii) δ ≅ E 1/2 – E– 1/2
interval (a, b)] (iii) δ = ∆E– 1/2 = E– 1/2 ∆
The bisection method calls for a repeated (iv) δ = ∇E1/2 = E1/2 ∇
halving of sub-intervals of [a, b] and locating the
half containing p. (v) δn f (x) = ∆n f x – nh
2
1
( )
Finite Difference Operators
(A) Finite difference (Forward operator)— = ∇n f x + nh
1
2( )
∆f (x) = f (x + h) – f (x) (vi) ∆∇ = ∇∆ = ∆ – ∇ = δ2
∆ 2 f (x) = f (x + 2h) – 2f (x + h) + f (x) (b) Averaging operator µ —
∆3 f (x) = f (x + 3h) – 3f (x + 2h)

n
+ 2f (x + h) – f (x) (i) µ f (x) =
1
2 [( 1
f x+ h +f x– h
2 ) (
1
2 )]
∆n f (x) = Σ (– 1) n – r n C r f (x + rh) (ii)
1 1/2
µ = [E + E– 1/2]
r=0 2
(B) The shift operator—
Interpolation and Extra Polation
Ef (x) = f (x + h)
En f (x) = f (x + nh) Interpolation means to find approximate value
of function f (x) for an x between different x values
(C) Relation between ∆ and E—
x 0, x 1 … x n at which the value of f (x) are given
1+∆ = E
i.e. f (x i) = f i (i = 0, …… n)
and ∆ = E–1
(1 + ∆)n = En If x 0 , x 1 , … x n are (n + 1) distinct values of
real valued function f ( x). We have a polynomial
and ∆n = (E – 1)n
p n (x i) ≈ f n (x) of degree n or less.
(D) Factorial notation—
i.e. there is most one polynomial of degree ≤
x (n) = x (x – h) (x – 2h)
n which interpolates, f (x) at (n + 1) distinct points

… (x – n – 1)h x 0 …… x n.
90 | Physical Sciences

Interpolation with Unequal Interval values and also extrapolating values of y a little
ahead (to the right) of yn.
(1) Divided difference (Aitken’s symbol)—
Extrapolation is the process of computing y- Let he function y = f (x) takes the values y 0 ,
outside the range a ≤ x ≤ a + nh i.e. when x < a y 1 , y 2 , …… corresponding to the value x 0 , x 0 + h,
and x > a + nh. x0 + 2h, … of x.
Various interpolation methods Suppose it required to evaluate f (x) for x = x 0
Interpolation + ph, where p is any real number. Then
y p = f (x n + ph)
= E p f (x n )
Forward Backward Central
difference difference difference = (1 – ∇)p y n (·. · E – 1 = 1 – ∇)

Newton
,
Newton s ,
Stirling s
= [1 + p∇ + p (p2!+ 1) ∇ 2

forward backward formula

formula
,
formula
,
Bessel s
p (p + 1) (p + 2) 2
+
3!
∇ + … yn ]
Lagrange s [by using Binomial theorem]
formula
formula for p (p + 1) 2
unequal interval y p = y n + p ∇y n + ∇ yn
2!
(1) Newten’s Forward Interpolation p (p + 1) (p + 2) 3
+ ∇ yn + …
Formula 3!
It is called Newton’s backward interpola-
This formula is used for interpolating the
tion formula as it contains y n and backward
values of y near the beginning of a set of tabulated
values and extrapolating values of y a little differences of y n.
backward (i.e., to the left) of y0. Stirling Formula
Let the function y = f(x) take the values y0, y 1 , By Newton’s formula
y2, … corresponding to the values, x 0 , x 0 + h, x0 + p (p – 1) 2
2h, … of x. Let it is required to evaluate f (x) for y p = y 0 + p.δy 1/2 + δ y1
2!
x = x 0 + ph where p is any real number
p (p – 1) (p – 2) 2
Ep = f (x) = f (x + ph) + δ y 3/2 + … …(1)
3!
yp = f(x0 + hp) = Ep f (x0) Since ∆y 0 = δy 1/2,
= (1 + ∆)p y0 + … [E = 1 + ∆] ∆ 2 y 0 = δ2 y 1
= {1 + p∆ + (p2!– 1) ∆ 2 1
Now δy 1/2 = (δy – 1/2 + δy 1/2)
2
+
p(p – 1) (p – 2) 2
3!
∆ + … y0 } 1
+ (δy 1/2 – δy – 1/2)
2
p(p – 1) 2 1
= y 0 + p∆y0 + ∆ y0 = µδy 0 + y 0
2! 2
+
p(p – 1) (p – 2) 3
∆ y0 + … δ2 y 1 = δ2 y0 + (δ3 y1 – δ2 y 0)
3! = δ2 y 0 + δ3 y 1/2
It is called Newton’s forward interpolation 1
formula as it contains y0 and the forward diffe- = δ2 y 0 + (δ3 y – 1/2 + δ3 y 1/2)
2
rence of y0.
1
(2) Newton’s Backward Interpolation + (δ3 y1/2 – δ3 y – 1/2)
2
Formula 1
= δ2 y0 + µδ3y0 + δ4y0
This formula is used for interpolating the 2
values of y near the end of a set of tabulated and so on.
 Physical Sciences | 91

From equation (i), we get To eliminate this formula, we shall first

p2 2 define these differences.
y p = y 0 + p µδy 0 + δ y 0
2! If (x 0, y 0), (x 1, y 1), (x 2, y 2)… be given points,
p (p2 – 12 ) 3 p2(p2 – 12 ) (p2 – 22 ) then first divided difference for the arguments
+ µδ y0 +
3! 5! x 0, x 1 is defined by the relation
× µδ5 y0 + … y – y0
[x 0 , x 1] = 1
Bessel’s Formula x1 – x0
By Newton’s formule y2 – y 1
p (p – 1) 2 Similarly, [x1 x2] =
y p = y 0 + p.δy 1/2 + δ y1 x2 – x1
2!
y3 – y2
p (p – 1) (p – 2) 3 and [x 2 , x 3] = etc.
+ δ y 3/2 + … x3 – x2
3!
1 Second divided differences for x 0 , x 1 , x 2 is
If δ2 y1 = µδ2 y 1/2 + δ 3y 1/2 defined as,
2
1 [x 1 ‚ x 2] – [x 0‚ x 1]
δ3 y 3/2 = δ3 y 1/2 + µδ4y1/2 + δ5y 1/2 etc. [x 0 , x 1, x 2] =
2 x2 – x0
Then on simplification, we get
The third divided difference for x 0 , x 1 , x 2 is
p (p – 1) 2
y p = y 0 + pδy 1/2 + µδ y1/2 defined as
2!
[x 1 ‚ x 2‚ x 3] – [x 0‚ x 1‚ x 2]

+
( )
p (p – 1) p –
1
2
δ 3y 1/2 +
p (p 2 – 12 ) (p – 2)
[x0 , x 1 , x 2 , x 3 ] =
x3 – x0
3! 4! so on.
µδ4 y1/2
Newton’s Divided Difference Formula

+
( )
p (p2 – 12 ) (p – 2) p –
1
2
δ5y1/2 + …
Let y 0 , y 1 , … y n be the values of y = f ( x)
5! corresponding to the arguments x 0 , x 1 , … x n .
Then from the definition of divided differences,
Interpolation with Unequal we have
Intervals y = f (x) = y 0 + (x 1 – x 0) [x 0 , x 1] + (x – x 0)
Lagrange’s Formula for Unequal (x – x 1) [x 0 , x 1, x 2] + (x – x 0) (x – x 1) (x – x 2)
Intervals [x 0 , x 1, x 2, x 3] + … + (x – x 0 ) (x – x 1)
If y = f (x) takes the values y 0 , y 1 , … y n … (x – x n ) [x 0 , x 1, … x n]
corresponding to x = x 0 , x 1, … x n, then
Numerical Integration
(x – x 1) (x – x 2) … (x – x n )
f (x) = y
(x 0 – x 1) (x 0 – x 2) … (x 0 – x n) 0 It is the numerical evaluation of integrals
b

+
(x – x 0) (x – x 2) … (x – x n )
y
I= ∫ a f (x) dx, where a and b are given and δ is a
(x1 – x0) (x 1 – x2) … (x 1 – x n) 1 function given analytically by a formula or
(x – x0) (x – x1) … (x – xn – 1) empirically by table of values i.e.
+…+ y
(xn – x0) (x n – x1) … (x n – x n – 1) n b

Divided Differences
J = ∫a f (x) dx
b
The Lagrange’s formula has the drawback
that if another interpolation value were inserted,
= ∫a φ (x) dx
then the interpolation coeff. are required to be where φ (x) is a polynomial which is equivalent to
recalculated. f (x).
92 | Physical Sciences

∫
x0 + hn
Rectangular rule—Subdivide the interval of
integration a ≤ x ≤ b into n-subintervals of equal = f (x)dx
x0
b–a 3h
length h = and in each subinterval approxi- = [(y 0 + y n) + 3(y 1 + y 2 + y 4 + y 5)
n 8
mate f by f (x j* ) = value of at the mid point x* j of … + y n – 1) + 2(y 3 + y 6 + y n – 3)]
the j-th subinterval Note—If polynomial is of degree one, then
b
J = ∫ a f (x)dx Trapezoidal rule is
1
J = h (y 0 + y 1)
= h [f(x 1* ) + f (x 2 * ) + … f (x n* )] 2
b–a If polynomial is of degree 2, then Simpson’s
where h = one-third rule gives
n
1
Trape Zoidal Rule J = h (y 0 + 4y 1 + y 2)
3
Subdivide the interval a ≤ x ≤ b into n- For polynomial of degree 3, then
subinterval of equal length 3
J = h [y 0 + 3y 1 + 3y 2 + y 3]
8
b–a,
h = Numerical Solution of Ordinary
n
Differential Equation
where x 0 = a and x 0 + nh = b
b
Range-Kutta Method—Given initial value
J = ∫a f (x)dx problem
y′ = f (x, y)

∫ ⇒
x0 + hn
= f (x)dx y (x 0) = y0
x0 Here k1 = hf (x n, y n)
=
h
2
[(y 0 + y n) ( 1
k 2 = hf x n + h ‚ y n + k 1
2
1
2 )
Simpson’s One-third Rule
+ 2(y 1 + y 2 + … + y n – 1)]
( 1
k 3 = hf x n + h‚ y n + k 2
2
1
2 )
k 4 = h (x n + h 1 y n + k 3)
Sub divide the interval a ≤ x ≤ b into even 1
number of equal intervals n = 2m of length and y n + 1 = y n + (k 1 + 2k 2 + 2k 3 + k 4)
6
b–a xn + x0
h = Where h = (i.e. x n = x 0 + nh)
2m n
b
∫
J = a f (x)dx
Elements of Probability and Theory of
Errors—A Priori Probability or Mathematical
x0 + nh Probability :—
= ∫x0
f (x)dx
P(A) = =
p No. of outcomes favourable to A
q total Number of outcomes
h A Posterior Probability (statistical def of
= [(y 0 + y n) + 4(y 1 + y 3
3 Probability)—If an experiment is performed N
+ …+ y n – 1) + 2(y 2 + y 4 + … + y n – 2)] times and the times that the experiment succeeds
is Ns, and if the ratio
Simpson’s Three-Eight Rule
N
Subdivide the interval a ≤ x ≤ b into multiple P(N) = s
N
of 3 of equal intervals n = 3m of length approaches the limit,
b–a N
h = lim P(N) = lim s = Ps
3m N→∞ N→∞ N
b
∫
J = a f (x)dx From def. of probability we conclude the
results.
 Physical Sciences | 93

(1) If ps is the probability of the success of an This f (x) is called probability function or
experiment, then (1 – p s) = p f probability distribution of the discrete random
(2) The value of prob. P is a number between variable.
0 and 1 i.e. 0 ≤ P ≤ 1 e.g. In an experiment, three coins are tossed
(3) The probability of a certain event is one— simultaneously. If the number of heads is the
random variable, find the probability function for
P(E) = 1 this random variable.
where P(E) is the probability of a certain S = HHH HHT HTH THH HTT THT TTH TTT
event, i.e. the event which occurs in every
3 2 2 2 1 1 1 0
observation
(4) The prob. of an impossible event is zero Since for 8 possible cases, the prob. of each
1
P(0) = 0 possible out come is
8
Compound Probability— 1
P(AB) = P(A) P (B/A) P(x = 0) = P(x 8) =
8
= P(B) P(A/B) P(x = 1) = P(x 5) + P(x 6) + P(x 7)
n (AB) 3
P(B/A) = =
n (A) 8
n (AB) 3
P(A/B) = P(x = 2) =
n (B) 8
Theorem of total Probability— 1
P(A + B) = P(A) + P(B) – P(AB) P(x = 3) =
8
Random Variables So discrete random variable
A function which can taken on any value x 0 1 2 3
from the sample space and its range is same set of 1 3 3 1
real number is called a random variable of the f (x)
8 8 8 8
experiment.
Random variable classified as CDF Commulative Distribution
(1) Discrete variable Function
(2) Continuous variable The Commulative Distribution Function
A Discrete random variable—May be (CDF) of a random variable ‘X’ may be defined as
defined as the random variable which can take on the prob. that a random variable ‘X’ takes a value
only finite number of values in a finite observation less than or equal to x.
interval. so CDF; F X(x) = P(X ≤ x)
e.g., experiment of tossing three coin simul-
taneously, here eight possible outcomes. Property
S HHH HHT HTH THH THT TTH TTT (1) 0 ≤ F x (x) ≤ 1
Continuous random variable—A random (2) F X (– ∞) = 0
variable that takes on an infinite number of values
is called a continuous random. F(∞) = 1
For x = – ∞, means no possible events. Due to
Probability Distribution of a Discrete this face P(X ≤ – ∞) will always be zero.
Random Variable x = ∞ means P(X ≤ ∞ ) which includes
Let x discrete random variable and also let probability of all possible events and probability
x 1, x 2, x 3 … be the values that X can be take. of a certain event is ‘1’ therefore
Then P(x = x j) = f (x j) j = 1, 2, 3, … F x (∞) = 1
will be the prob. of x j. (3) F X (x 1) ≤ F X (x 2) if x 1 ≤ x 2
94 | Physical Sciences

CDF or Discrete Random Variable The Bernoulli distribution is given by

For discrete random variable, the cumulative f (x, p) = p x (1 – p) 1 – x
Distribution function (CDF) for a complete range = p x q 1 – x, x = 0, 1
of x may be expressed as Mean µ = p
F X (x) = 0 for – ∞ ≤ x < x 1
n
and Variance, σ2 = pq
= Σ p (X = x j) for x ≤ x ≤ x n Binomial distribution—A random variable
j=1
X has a Binomial distribution and is a Binomial
= 1 for x n < x < ∞ random variable if and only if its distribution is
Probability density function— given by
d
PDF; f X (x) = F (x)
dx X f (x, n, p) = (nx) p qx n–x

PDF is more convenient representation for or n C x p x qn – x ; x = 0, 1, … n

continuous random variable.
with p+q = 1
Properties
The number of success in n trials is a random
(1) Prob. density function, is always non zero variable, having Binomial distributions with para-
for all values of x meters p and n. The term f (x, n, p) is a successive
f X (x) ≥ 0 for all values of x. term in Binomial expansion [p + q] n
(2) The area under the PDF curve is always Mean µ = np
equal to unity Variance σ2 = npq
∫
∞
f (x)
–∞ X
dx = 1 Poisson distribution—The poisson
distribution is given by
(3) The CDF (comulative distribution
e– λ λx ,
function) may be obtained by integrating PDF f (x; λ) = x = 0, 1, 2, …
x!
∫
x
F X (x) = f
–∞ X
(x) dx The mean µ = λ

∫ σ2 = λ
x2
Variance
(4) P(x1 < X ≤ x 2 ) = x1
f X (x) dx
Normal distribution—The probability
density function for normal distribution is
Probability Distribution
Uniform distribution—A random variable
that takes on k-different values with equal proba-
f (x) =
σ
1
√2π
exp [– 12 (x σ– µ) ]
2

bility, then it has a uniform (discrete) distribution. (σ > 0)

The discrete uniform probability, distribution where
is given by (1) µ is the mean and σ the standard
1 deviation.
f (x) =
k 1
x = x 1, x 2, … x k, (2) is a constant factor that makes the
σ
√2π
xi ≠ xj i≠j area under the curve equal to 1
k
xi
∫
∞
and have mean, µ = Σ i.e. f (x) dx = 1
i=1 k –∞
k
(x i – µ)2 (3) The curve of f (x) is symmetric with
and variance, σ2 = Σ
i=1 k respect to x = µ. For x = 0 = µ, it is symmetric
Bernoulli distribution—If an experiment has with respect to y-axis, x = 0 (bell shaped)
two possible outcomes success and failure with (4) The exponential function tend to zero very
probabilities p & q = 1 – p, then the number of fast, the faster the exponential function, smaller
successes 0 or 1 has a Bernoulli distribution. the standard deviation σ is.
 Physical Sciences | 95

The probability distribution functions The curve φ(z) is s spaced, increase monotone

∫ 1 u–µ 2
∞ 1
F(x) =
σ
1
√2π – ∞
exp.[ ( )]
–
2 σ
du from 0 to 1 and intersect the vertical axis at ·
2
For standard normal distribution (µ = 0 and
σ = 1)
1
The distribution function
φ (z) =
1 ∞
∫
2π – ∞ ( )
exp –
u2
2
du 1
2

f (x) =
1
2π
exp ( )
–x
2
−2 O 2

Solved Questions
→ Solution : Let
Q. 1. If V = Vr ^er + Vθ e^θ, where Vθ and
→
Vr are the functions of t and A = Ax^i + Ay ^j + Az^k
^
e = ^i cos θ + ^j sin θ →
r B = B x^i + By^j + Bz^k
^e = – ^i sin θ + ^j cos θ →
θ
C = C x^i + Cy^j + Cz^k
dV
find · Then
dt
→ → → →
Solution : Given that (A × B ) · (C × A )
→ = [(AyBz – AzB y) ^i + (AzB x – BzAx) j^
V = Vr e^r + Vθ e^θ
→ + (A xBy – AyBx) ^k] [(CyAz – AyCz) ^i
d V ^ dVr de^ dVθ de^
So = er + Vr r + e^θ + Vθ θ …(1)
dt dt dt dt dt + (C A – A C ) ^j + (C A – C A ) ^k]
∴ z x z x x y y x
e^r = ^i cos θ + ^j sin θ = [(AyBz – AzB y) (CyAz – AyCz)
e^ = – ^i sin θ + ^j cos θ
θ + (A zB x – BzAx)(CzAx – AzC x)
de^r dθ dθ + (AxBy – AyBx)(CxAy – CyAx)]
∴ = – ^i sin θ + ^j cos θ = AyBzC yAz – Ay2B zC z – Az2 B yCy + AzAyByCz
dt dt dt
^ dθ + AzB xCzAx – Az2 B xCx – BzAx2C z + BzAxAzC x
= – eθ
dt + AxByCxAy – Ay2B xCx – Ax2B yCy + AyBxCyAx
^
deθ dθ dθ
and = – ^i cos θ – ^j sin θ = (AxBx + AyBy + AzB z) (AxCx + AyCy + AzC z)
dt dt dt – (BxCx + ByCy + BzC z) (Ax2 + Ay2 + Az2 )
dθ → → → → → → → →
= e^r = (A · B ) ( A · C ) – ( B · C ) ( A · A )
dt
Substituting these values in (1), we get Q. 3. A curve, C is defined by the para-
→ metric equations x = x (S), y = y (S) and z = z (S)
d V ^ dVr ^ dθ dVθ ^ dθ
= er + eθ Ar + e^θ – er Vθ where S is the length of arc measured from a
dt dt dt dt dt fixed point on the curve C. If r is the position
Q. 2. Prove that vector of any point on the curve, show that
→ → → → → → → → →
(A × B ) · ( C × A ) = (A · C ) ( B · A ) d r /dS, is an unit vector, tangential to C, at the
→ → → → →
– (B · C ) (A · A ) point r , i.e., (x, y, z).
 Unit-2

CLASSICAL DYNAMICS
The branch of physics in which mass, length y
and time are taken as absolute quantities is known
as classical mechanics. Classical mechanics is eθ
er
developed on famous Newton’s laws of motion.
Motion of a Particle
(a) Velocity—The instantaneous velocity of a P
particle is given by a vector
→ r
dr
v =
dt θ
δr x
= limδt → 0 O M
δt
(r + δr) – r The above expressions for ^e r and ^e θ follow
= limδt → 0
δt from the transformations,
→ x = r cos θ,
where r is the position vector of the particle at an
instant of time t. Components of velocity in y = r sin θ
rectangular co-ordinates
→
ṙ = ẋ ^i + ẏ ^j + ż ^k, which implies that ẋ, ẏ, ż Now, r = r e^r
represent the speed of the particle parallel to
Thus, ṙ = ṙ er + r ėr
respective co-ordinate axes.
Component of Velocity in Plane Polar = r (eθ θ̇) + ṙ er = v
Coordinates Thus, the components of v are radial velocity
Here, (x, y) ≡ (γ, θ) in x – y plane such that : (vr) = ṙ along e^r, in direction of r increasing.
→
^e r = Unit vector along OP Transverse velocity (v θ ) = rθ̇ along ^e θ , in
 ∂→  direction of θ increasing.
 r
 ∂r  Components of velocity in spherical polar
=
 → co-ordinates.
∂ r  Here, (x, y, z) ≡ (r, θ, φ) such that :
∂r
x = r sin θ cos φ
= cos θ^i + sin θ^j
→ y = r sin θ sin φ 

and ^e θ = Unit vector along PM z = r cos θ
 ∂→  Above set of equation suggests
 r
=
 ∂θ  r = (x2 + y2 + z2 )1/2,
φ = tan–1 (y/x)

 →
∂ r  √ y + x   r = x i + y j + zk
^ ^ ^
  
2 2
∂θ and θ = tan–1
 z 
= – sin θ^i + cos θ ^j . 
152 | Physical Sciences

 ∂→ 
 r
 ∂r 
and φ = tan–1 (yx)
^e r =
 →
∂ r 
∂r
= sin θ cos φ^i + sin θ sin φ^j + cos θ ^k
= Unit vector along r-increasing,

 ∂→ 
 r
^e ρ =  ∂ρ 
 →
∂ r 
∂ρ
 ∂→ 
 r = cos φ^i + sin φ^j
^e = ∂θ 

θ
 → = Unit vector along ρ increasing,
∂ r   ∂→ 
∂θ
 r
= cos θ cos φ^i + cos θ sin φ^j – sin θk^ ^e φ =  ∂φ 
= Unit vector along θ increasing  →
 ∂→  ∂ r 
∂φ
 r
 ∂φ 
and ^e φ = = – sin φ^i + cos φ ^j
 →
∂ r  = Unit vector along ρ increasing
∂φ
= – sin φ^i + cos φ ^j
 ∂→ 
 r
= Unit vector along φ increasing ^e z =  ∂z 
and
→  →
Now, r = r^e r thus ṙ = ^e r ṙ + r^ė r ∂ r 
∂z
⇒ v = ṙ
= ^k
= ṙ ^e r + rθ̇ ^e θ + r sin θ φ̇ ^e φ
= Unit vector along z-increasing
Components of velocity in cylindrical co-
ordinates. Thus, v = ṙ
Here (x, y, z) ≡ (ρ, φ, z), so that = ρ̇ ^e ρ + ρ̇ φ e^φ + zk^ [Q e^z = k^ ]
x = ρ cos φ 
 (b) Acceleration—Acceleration in Cartesian
y = ρ sin φ 
 coordinates
z = z
→
Such that v̇ = a
ρ2 = x2 + y2 x ^i + ˙˙
= ˙˙ y ^j + ˙˙
z ^k
 Physical Sciences | 153

Acceleration in plane polar co-ordinates Thus, the rate of change of angular momentum
→ of a particle equals the torque acting on it.
v̇ = a
r – r θ̇2 – r sin θ φ̇) ^e r + (r˙θ̇ + 2 ṙ θ̇
= (˙˙ Work done by the force acting on a particle
→
– r sin θ cos φ2) e^ + (r sin θ ˙˙
θ φ + 2 ṙ φ̇ sin θ and kinetic energy—If the force F acting on a
→
+ 2 r cos θ φ̇ θ̇) ^e φ particle displaces it through a distance dr then the
= R e^r + S e^θ + T e^φ work done by the force on the particle is defined
as
Acceleration in cylindrical co-ordinates.
→ → →
a = (˙ρ̇ – ρφ̇2) e^ρ + (ρ˙˙
φ + 2 ρ̇φ̇) e^φ + ˙˙
z ^k dW = F . dr .
Linear Momentum of a Particle The work done by a force F in moving the
→ → particle from A to B along a curve C is given by
The vector p = m v is known as linear the time integral.
momentum of the particle. Applying Newton’s
∫
B
law of motion, we get WAB = A
F . dr
→
d →
∫
d (m v ) rB
(p ) = = F . dr
dt dt rA
= mv̇ where rA and rB are the position vectors of A
→ and B.
= ma = F ext
→ If these points corresponds to velocities v A
here F ext is the external force acting on the and v B of the particle at the instants t A and tB
particle. Moment of force or torque and angular respectively, then
momentum of the particle. If Fext is the external vB
force acting on the particle then its moment
→ → →
WAB = ∫ vA
m (dvdt) v dt
τ = r × F is known as torque. vB

∫
1 d 2
→ d → = m (v ) dt
τ = r × (m v ) 2 vA dt
dt
1
The direction of torque is perpendicular to the
→ →
=
2 (
m v 2B – v 2A )
plane containing r and F . = TB – TA
The moment of momentum or angular
momentum is defined as Change in kinetic energy.
→ → → Power
L = r ×p
→ → The time rate of doing work on a particle is
= r ×mv
called the instantaneous power or power and is
r is the radius vector from the fixed point to given by
the particle.
dW d
d → d → → P = = (F . dr)
Now, (L) = (r ×p) dt dt
dt dt
→
→ dp → d r
= r × + p ×
→ = F.
dr
dt
=F.v ()
dt dt Impulse
→ → tB
= r × F +0
→ →
The time integral ∫ tA
F.dt is called the
= τ [Q u × m v = 0] impulse of force F. Here
154 | Physical Sciences
tB tB

∫ ∫ ∫
d B
F.dt = (mv) dt) WAB = F.dr
tA tA dt A
v
[mv] vB
∫
= B
A = – ∇V . dr
= mv B – mv A = pB – pA A

∫ (∂V∂x ^i + ∂V∂y ^j + ∂V∂z k^)

Change is momentum. B
= –
This is true even when the mass varies and A
the force is non-conservative.
(dx ^i + dy ^j + dz k)
^
Force Field and Conservative Force
∫ (∂V∂x dx + ∂y∂Vdy + ∂z∂Vdz)
B
A force acting on a particle which at every = –
A
point of the continuum is a function of the
∫
B
position r = (x, y, z) is said to constitute a force
= – dV = VA – VB
field. A
In general the value of line integral for the Thus, work done is independent of the path
∫
B
work done WAB = F.dr depends on the path of chosen.
A Conservation theorems for a particle :
integration. On the contrary a force for which the (a) If the net external force F acting on a
work done in displacing a particle from one point particle is zero, its linear momentum P does not
to another is independent of the path followed, is change with time and is conserved.
called a conservative force; otherwise, the same is
non-conservative. Hence, the attention should be (b) If the net external torque τ acting on a
paid to the following theorems : particle is zero, then its angular momentum L does
not change with time and is conserved.
(a) If F = –∇V, where V is single valued (c) As we have seen
and has continuous partial derivatives, then the
work done in moving a particle from one point WAB = TB – TA = VA – VB
A(xA, yA, zA) in this field to another point P(x, y, z) = TA + VA = TB + VB .
is independent of the path joining A, P. The scalar quantity E = T + V, i.e. sum of
Conversely, if ∫ F.dr is independent of the path C kinetic and potential energies is called the total
joining any two points then there exists a function energy of particle. Thus, if the forces acting on a
V such that F = –∇V. particle are conservative, its total energy remains
P2 conserved.
(b) If∫ P1
F . dr is independent of the path The equation of motion of a particle
“D’Alembert’s principle”
joining any two points P1 and P2 in a given region,
∫
By Newton’s second law of motion the rate of
then O F.dr = 0; for all paths in the region and dp
change of momentum and the impressed force
conversely. dt
(c) A necessary and sufficient condition for F have the same direction and their magnitudes
F 1 dx + F2 dy + F3 dz to be an exact differential is are proportional. Thus,
→ →
that dp d →
→ → dt
=
dt
(m v ) = k F
∇ × F = 0 where F = (F 1 , F 2 , F 3 ).
dv
Potential Energy or m = kF
dt
The scalar V such that F = –∇V is called the → →
potential energy or scalar potential or simply or m f = kF ,
potential of the particle in the conservative force where k is a scalar constant. In many system of
field F. In such a case the total work done in units, we choose the unit of force such that when
moving the particle along path C from points A to it acts on a particle of unit mass it will produce
B is given by unit acceleration. Here let k = 1.
 Physical Sciences | 155

→ → not changed dz = 0 which gives a constant on

Thus, mf = F integration) or holonomic constraint. On the
→ → contrary there also exists a non-holonomic or non-
∴ F –mf = 0 integrable constraints viz., motion of a particle on
i.e., the external force on a particle along with the sphere. The mechanical system in which the
→ time t does not enter explicitly is called to have a
the reversed effective force (–m f ) form a system scleronomic constraint on the other hand the
in equilibrium. This is called D’alembert’s
system in which the moving constraints are
principle for a particle.
involved and the time t enters explicitly is called
Motion of a particle under time dependent to have a rheonomic constraints.
applied force—By Newton’s second law, we (ii) Generalised co-ordinates—If a particle or
have a system of N particles moves subject to possible
mdv = Fdt. constraints, then there must be a minimum number
Let v0 be the initial velocity of the particle at of independent co-ordinates needed to specify the
time t0 and v be its velocity at any time t, then on motion. These co-ordinates denoted by q1 , q2 , …,
integrating the above equation qn are called generalised co-ordinates and may be
distances, angles or quantities related to them.
∫ ∫
v t
v0 dv = t0 F(t)dt (iii) Degrees of freedom—Consider a system
of N particles. We generalise, first of all, to
∫
t
or mv – mv0 = t0 F(t)dt generalised co-ordinates, for the system, using
the notation, x1, x2 x 3 for the space co-ordinates
dr (x1 , y1, z1 ) (In Cartesian system) of the first
Since, v = ,
dt particle, x4, x5, x6 for the space co-ordinates (x2, y2,
∫ F(t)dt
t
dr 1 z2 ) of the second particle and so on, moving to
Thus v = = v0 + t0
dt m x 3N – 2, x 3N – 1, x3N for the space co-ordinates (x n ,
dt + [ ∫ F(t)dt ]dt
t
1 yn, zn ) of the Nth particle. These are thus 3N space
∴ dr = v0 t0
m co-ordinates for the N particles system, given by
Incase r0 be the position vector of the particle x1, x2, x3 …, x3N we also use the notation m1, m2,
at time t0 and r be the same at time t, then m3 …, m N for the masses of the N particles
respectively.
∫ ∫ ∫ [∫
r t t t
r0 dr = t0 v0 dt +
1
m t0 t0 F(t)dt dt ] The minimum number of independent vari-
ables (Not necessarily the co-ordinates) required
1 t
∫ [∫
t
∴ r = r0 + v0 (t – t0) +
m t0 t0 F(t)dt dt ] to fixe the position and the configuration of a
dynamical system which are compatible (i.e., obey
The above equations for v and r gives us the the restrictions) with the given constraints is
velocity and position vectors of the particle at any called the number of degrees of freedom of the
instant of time. system. Here the term ‘independent’ is taken to
mean linear independence (i.e.). If the variables be
Some Definitions denoted by q 1 , q 2 , … qr, … qn where n is the
(i) Constraints—The motion of a free number of degrees of freedom, then it is possible
particle in space is described by three independent to find a set of non-vanishing numbers c1 , c2 … cn,
co-ordinates and it is said that the particle has three real or complex for which the sum (c1q1 + c 2 q2
degrees of freedom. When the particle is not + … cnqn ) (i.e.) ∑ crqr should be identically zero.
allowed to move freely in three dimensions, we say r
that it is subjected to constraints and thus number If there be k equations of constraints, then for
of degrees of freedom is reduced. When a particle the total of 3N co-ordinates (x1, x2, … x 3N ) to be
motion is confined to a plane only. It has said to completely determined. It is enough if we consider
have two degrees of freedom instead of three. In n = (3N – k) variables, q1 , q2 … q n in terms of
general when a particle is to move on a surface which the 3N co-ordinates may be expressed as
only, i.e., f (x, y, z) = 0, then such constraints are transformations. Thus, we express the space
said to be an integrable constraints (e.g. if z does co-ordinates,
156 | Physical Sciences


x = x (q , q , q , …, q , t)
x1 = x1(q1, q2 , q3 , …, qn , t) The kinetic energy of the system can also be
written as a quadratic form in the generalised

…………………………..… 
2 2 1 2 3 n
velocities qα. If the system is independent of time

…………………………..…. 
explicitly, i.e., scleronomic, then the quadratic

…(1)
form has only terms of the type aαβ qα qβ. For the

x = x (q , q , …, q , t)  
x = x (q , q , q , …, q , t)
i i 1 2 3 n
rheonomic system, linear terms in q α are also
present.
3N 3N 1 2 n
Generalised Forces
In general, the xi’s are explicit function of
If W is the total work done on a system of
time, then the system is rheonomic. If all of the particles by force Fi acting on the ith particle, then
above xi’s do not solve time explicitly, the system N
is Scleronomic. The reverse transformations of set dW = ∑ φα dqα
α=1
(1) is given by N ∂r
q1 = q1 (x1, x2, x3, …, x3N, t)
 where φα = ∑ Fi ∂qi
i=1 α


q2 = q2 (x1, x2, x3, …, x3N, t) is called the generalised force associated with
q3 = q3 (x1, x2, x3, …, x3N, t) generalised co-ordinates qα. The generalised coor-


……………………………. …(2) dinates and generalised force need not have the


……………………………. dimensions of length and force respectively but
qn = qn (x1, x2, x3, …, x3N, t) the scalar product of the two will have the
dimensions of work. In the present discussion α is
the number of degrees of freedom and i is the
These are n equations in number and when number of particles in the system.
the q i’s are known, are not alone sufficient to Generalised Momentum
determine the 3N co-ordinates from (2). However, ∂T
there are k equations of co-ordinates which are We define p α = as the generalised
∂q̇α
also relations in terms of x i’s. Thus, there are
momentum associated with the generalised
totally n + k = (3N) equations which, therefore, co-ordinates q α. We usually refer p α as the
are sufficient to determine the 3N co-ordinates at conjugate momentum to qα.
any time t.
The n variables q1 , q2 , q3 , …, qn are called the
Frame of Reference
generalised co-ordinates of the system. A judi- A frame of reference is a real rigid body with
cious choice of the n variables, will enable the respect to which the motion of the body being
equations of motion in terms of the q i’s to be studied is to be considered. Rigidly fixed in the
frame of reference is some kind of co-ordinate
integrated easily to obtain these qi’s. The xi’s may system, so that the position of each point of a
then be determined from equations (2) and the moving body can be uniquely determined by the
equations of constraints. three co-ordinates of the point. Moreover the
Transformation Equations frame of reference should be furnished with a
‘clock’ by means of which the instants of time
The relationship of the generalised co-ordi- corresponding to various positions of a moving
nates q1 , q2 , q3 … qn to the position co-ordinates body in space are uniquely determined. Those
are given by the transformation equations. frames in which Newton’s law of inertia hold true
is called an inertial frames of reference or
Kinetic Energy and Generalised Newtonian frames of reference.
Velocity
Now let us investigate the case of moving
The total kinetic energy of the system is frames of reference under the following heads—
1 N (a) Frames of reference with uniform
2 i∑
T = mi ṙ 2i .
=1 translation velocity—Let S be a Newtonian
 Physical Sciences | 157

frame of reference and S′ another frame of →

reference in uniform translation motion relative Now differentiating P with respect to time
to S. The law of motion in those frames yields ‘t’, we get
→ → → →
F = m a = m a′ showing that the law of motion ^ ^ ^
d P dP1 ^ dP 2 ^ dP 3 ^ di dj dk
also helds in S′. Thus given one Newtonian frame = i+ j+ k + P1 + P2 + P3
of reference S, we can find an infinity of other dt dt dt dt dt dt dt
Newtonian frames of reference which have a →
uniform motion of translation relative to S. But i is a vector fixed in a rigid body S′, we
→
(b) Frames of reference with translation may take i as the position vector of a particle B
motion under constant acceleration—Let S be a of this body relative to a base point A (the origin
Newtonian frame of reference and S′ be another → di
of i ). So is the velocity of B relative to A and
frame of reference which is in relative translation dt
with respect to S. Then for a moving particle, we Hence, it can be written as
have
d^i →
→ → → = Ω × ^i
a = a 0 + a′ , dt
where a0 is the acceleration of S′ relative to S. d^j →
Now if m is the mass of the particle and F is the Similarly = Ω × ^j
dt
force acting on it, the law of motion is
→ → d^k → ^
F = ma and = Ω ×k
dt
(since S is a Newtonian frame)
Substituting these values in the above
The motion of particle with respect to S′ is equation, we obtain,
given by
→ → → →
a′ = m ( a – a 0 )
→ →
dP
dt
= (dPdt ^i + dPdt ^j + dPdt k^)
1 2 3

= ma – ma 0 → → →
→ → + P1 Ω × ^i + P2 Ω × ^j + P3 Ω × k^
= F – m a 0, →
which shows that the motion of S′ gives rise to the δP ^ ^ ^
→ = + Ω (P1 i + P2 j + P3k)
δt
following force – m a 0 . Hence, we can regard S′
as Newtonian frame provided we add to the actual δ
where we use to indicate partial differentiations
→ δt
forces a fictious force – m a 0 on each particle. ^ ^ ^
in which i , j , k are fixed.
(c) Frames of reference rotating with
→ →
δP → →
constant angular velocity—Let S be a Newtonian
dP
frame and S′ another frame of reference rotating Thus = +Ω ×P
dt δt
about a point θ of S with constant angular velocity
Ω. Let ^i , ^j , k^ be a triad of unit orthogonal vectors Obviously
dP
consist of two parts :
in S′. dt
→ δP
(i) is the rate of change (usually called the
Any vector P (e.g. displacement, velocity or δt
acceleration) can be represented in frame S′ as →
→ rate of growth) of P as measured by an observer
P = P 1^i + P2^j + P3^k moving with S′ keeping ^i , ^j , ^k as constants.
→ → →
we want to find the rate of change of P as viewed (ii) ( Ω × P ) is due to rotation of triad ^i , ^j , ^k
by an observer in S. and is called the rate of transport. Now, we can
158 | Physical Sciences

describe fully the motion of a particle relative to v'

the rotating frame as follows : Centrifugal
force
→ A
Let S′ rotate with an angular velocity Ω Coriolis
about a point O fixed in Newtonian frame S. force
r
Hence, the velocity of the moving particle A by
inspection from above is given by
→ → O
dr δr → →
v = = +Ω × r
dt δt Thus, we find that the rotation of S′ gives rise
→ → to two fictious forces –mat and –ma0, when the
(where O A = r ) two forces are added to true force F, the law of
Also acceleration be motion of a particle in S′ is precisely the
→ Newtonian law.
→ dv δ r → →
a = = +Ω × v In case the point O is also moving with
dt δt
→
→ →
δ  δ r → → → δ r → →
velocity v 0 and acceleration a 0 , then relative to a
Ω × Ω × +Ω ×r frame S′ the motion of the particle is governed by
δt  δt 
= + r +
 δt 
the general expression
→ → → δ→
δ2 r δ r → r → → → → →
= + × r + 2 Ω × m a ′ = F – ma0 – m at – ma0 .
δt2 δt δt
→ → → Lagrangian Dynamics
+ Ω × (Ω × r )
→ → → In Lagrangian formulation, it is assumed that
= a′ + a t + a 0 the mechanical state of a system is completely
δr defined, once its velocities and generalised
where v′ =
δt co-ordinates are specified.
which is the velocity of a particle with respect to S Lagrange’s Equation
→
δ2 r Let Fk be the net external force acting on kth
and a′ =
δt2 particle of a system, then from Newton’s second
is the acceleration of the particle with respect to S. law, we have
Here acceleration of the particle with respect mk ˙˙
r k = Fk
to S is given as
∂rk ∂r
→ → → or mk ˙˙
rk = Fα k …(1)
a′ = a – a k – a 0 ∂qα ∂qα
Thus, the force on the particle with respect to We also have
S is
→ → rk = rk (q1 , q2 , q3 … qn , t)
→ → → →
ma ′ = F – ma t – ma 0 (Q F = m a )
∂rk ∂r ∂r ∂r
→ ṙk = q̇ + k q̇ + … + k q̇n + k
∂q1 1 ∂q2 2 ∂qn ∂t
where m a ′ is called the force of transport and is
intimately related to centrifugal force given by ∂ṙk ∂rk
→ → → → → Thus = …(2)
( )
m Ω × Ω × r . The force ma0 = 2m Ω × v ′ is ∂qk ∂qα
called Coriolis force being perpendicular to both Now
→ → ∂ṙk ∂  ∂rk ∂r ∂r ∂r 
Ω and v ′. The figure below depicts the direction
∂qα  ∂q1 1 ∂q2 2 
= q̇ + k q̇ + … + k q̇n + k
of forces : ∂qα ∂qn ∂t 
 Physical Sciences | 159

∂2 r k ∂2 r k Hence, from equation (3)

= q̇1 + q̇ + …
∂qα∂q1 ∂qα∂q2 2 d ∂T ∂T ∂r ∂W
+
∂2 r k
q̇ +
∂
(∂r∂t )
k ( )
dt ∂qα
– =∑F k=
∂qα k k ∂qα ∂qα
= φα …(4)
∂qα∂qn n ∂qα
For α = 1, 2, …, n, we have n different
∂  ∂rk  ∂  ∂rk  equations which are called as Lagrange’s equa-
∂q1  ∂qα 1 ∂q2  ∂qα 2
= q̇ + q̇ + … tions. If the forces are derivable from a potential
V, then
∂  ∂rk  ∂  ∂rk  ∂W ∂V
∂qn  ∂qα n ∂t  ∂qα
+ q̇ + φα = =–
∂qα ∂qα

or
∂
∂qα (∂r∂t ) = dtd ∂q∂r 
k k
α
Since the potential is a function of q’s only
(and possibly the time t), then
Which shows after interchange of the order of
operator (∂q∂T ) – ∂∂Tq̇ = – ∂q∂V
d
dt α α α
∂ ∂ ∂ d ∂T ∂T ∂V
d
( )
dt ∂qα
=
∂qα ∂t () or
dt [∂q
(T – V)] – (
α ∂q ∂q )
–
α
= 0
α
Now,
d  ∂rk  ∂r d  ∂r  [Q ∂q∂V = 0]α
dt  k ∂qα
ṙ r k k + ṙk  k 
= ˙˙
∂qα dt  ∂qα d  ∂L  ∂L
dt  ∂q· α ∂qα
or – = 0 …(5)
∂r ∂ ∂rk
r k k + ṙk
= ˙˙
∂qα ∂qα ( )
∂t where L = T–V
∂rk d  ∂rk  ˙˙ ∂ṙk The function L defined as L = T – V is said to
dt  k ∂qα 
∴ ˙˙
rk = ṙ – rk
∂qα ∂qα be Lagrangian function—
(i) Lagrange’s equations for a non-holonomic
∂rk d  ∂r  ∂ṙ
or mk˙˙
rk = mk  ṙk k  – mk ṙk k systems with moving constraints are given by
∂qα dt  ∂qα ∂qα
Thus, from equation (1)
d
dt (∂∂Tq̇ ) – ∂q∂T
α α
= φα + λ1Aα + λ2B α + … …(6)

d  ∂rk  ∂r ∂r where α = 1, 2, …, n and λ1, λ 2 … are lagrange’s

dt  k k ∂qα
m ṙ – mk ṙk k = Fk k
∂qα ∂qα multipliers.
If the forces are derivable from potential, i.e.
Summing both sides with respect to k over all
∂V
particles, we have φα = –
∂qα
d  n ∂r  ∂ṙ ∂r
dt k∑
 mk ṙk k  – ∑ mk ṙk k = ∑ Fk k …(3) where V does not depend on q̇α, then equation (6)
=1 ∂qα k ∂qα k ∂qα
can be written as
Now, let T be the kinetic energy of a system d ∂L ∂L
of particles, then ( )
dt ∂q̇α
–
∂qα
= λ1Aα + λ2B α + …… …(7)
1 1
T = ∑ mk ṙ2k = Σ mk ṙk . ṙk (ii) When the system is conservative, the
2 k 2
generalised forces can be derived from a scalar
∂T ∂r potential function independent of velocities, the
∴ = ∑ mk ṙk k
∂qα k ∂qα equation of motion can be expressed in Lagrangian
∂r form. For a non-conservative systems in which the
∂T
and = ∑ mk ˙˙
rk k potential depends on velocities. The equations of
∂qα k ∂q α motion can also be expressed in Lagrangian form
[using equation (2)] provided the generalised forces φα are given as
160 | Physical Sciences

∂V d  ∂L  N ∂L
φα = – ∑ (pα dqα + pα dq̇α) + ∂t dt
∂qα dt  ∂q· 
+ =
α=1
α
N N
The quantity V is known as generalised
potential or the velocity dependent potential.
= d ( ∑ p q̇ ) + ∑ (ṗ dq – q̇ dp )
α=1
α α
α=1
α α α α

(iii) When the Lagrangian does not involve a ∂L

+ dt
certain co-ordinates there exists another kind of ∂t
Integral which is called as momentum integrals. N N ∂L
The co-ordinate which is missing is termed as
cyclic or ignorable co-ordinate.
or d
{ ∑ pαq̇ α – L
α=1 } = ∑ (q̇ α dpα – p α dqα) –
α=1 ∂t
dt

N
When L do not contain qα explicity, then
∂L
Thus, if the quantity [ ∑ (p q̇ – L)],
α=1
α α

= 0 when expressed in terms of (q1 , q2 , q3 , …, qn ), (p1 ,

∂qα p2 , …, pn , t) be denoted by H, we have
d ∂L
or ( )
dt ∂q̇α
= 0 (from Lagrange’s equations) dH =
N ∂L
∑ (q̇α dpα – pα dqα) – ∂t dt
α=1
∂L
or = constant = pα (say) N ∂H ∂H N ∂H
∂q̇α = ∑ dqα + ∑ dpα + dt
α = 1 ∂qα α = 1 ∂pα ∂t
Hence, if a co-ordinate q α is cyclic, the N ∂L
quantity
∂L
represents an integral of motion and is = ∑ (q̇α dpα – ṗα dqα) – ∂t dt
α=1
∂q̇α
denoted by p α . The constant p α is known as Thus
conjugate momentum to qα. ∂H ∂H ∂H ∂L
q̇α = , ṗ = – , =–
∂pα α ∂qα ∂t ∂t
Hamiltonian Formulation
These equations are called Hamilton’s equa-
Another mode of the description of the tions or Hamilton’s canonical equations and the
mechanical state of a system (Unlike Lagrangian function H is called Hamiltonian. The total of
which is described in terms of velocities and Hamilton’s equation is the same as the total order
generalised co-ordinates) in terms of momenta and of Lagrange’s equations viz., 2n. But whereas
generalised co-ordinates is much more advanta- Lagrange’s equations give us n equations each of
geous. Before proceeding further, we write second order, Hamilton’s equations are 2n each of
Lagrangian of the system as first order.
L = L (qα, q̇α, t) (i) If the Hamiltonian H is independent of t,
Also Lagrangian equations of motion are we have
d ∂L ∂L δH N ∂H dqα N ∂H dp
( )
dt ∂q̇α
–
∂qα
= 0
dt
= ∑
α = 1 ∂qα dt
+ ∑
α = 1 ∂pα dt
α

∂L N N
Writing p = ( ), we get
α
= ∑ ṗα q̇α + α∑= 1 q̇α ṗα = 0
∂q̇ α α=1

∂L
p = ( )  ∂H = – ṗα
α
∂q  ∂qα
α 
Hence, from these sets of equations, we can  ∂qα = – q̇α
regard either of the sets of quantities (q̇1 , q̇2 , q̇3 , ∂p
…, q̇n ) or (p1 , p2 , p3 , …, pn ) as function of the or H = constant.
other set. (ii) By Euler’s theorem on homogeneous
Now, ∂T
functions, we have Σ q α = 2T, where T is the
N
∂L ∂L ∂L ∂q̇α
L = ∑ [
α = 1 ∂qα
dqα +
∂q̇α
dq̇α +
∂t]dt
kinetic energy of the system.

∂L ∂T
But L = T–V ∴ =
∂q̇α ∂qα
(Q V does not depend on q̇α)
∂L
or Σ q̇α = 2T
∂q̇α
or Σ q̇α pα = 2T
∴ H = Σ ṗα q̇α – L = 2T – (T – V)
= T+V=E
Hence, the Hamiltonian H when independent
of t is equal to the total energy of the system.
Collisions and Scattering
Scattering is the phenomenon where an
uniform beam of particles having same mass and
energy incident upon a centre of force localised at
a point or in a very small region, get scattered in
different directions around the scattering centre.
Scattering, therefore, differs from collision. A
collision involves a couple of particles or a few of
them, or a couple of rigid bodies, each of which
may be separately considered. In a scattering
process, on the other hand, no isolated particle is
or can be considered separately. When a beam of
particles is scattered, we assume a sufficiently
large number of them, to talk of the intensity of
the beam not only incident but also scattered in
different directions. The scattering centre, assumed
at rest, in the laboratory system of co-ordinates,
also is assumed to be of finite extension even if a
central force towards a point the scattering centre
is assumed.
Scattering problems assumed importance
because of nuclear physics, the earliest scattering
experiment being discussed by Rutherford when
he allowed a beam of α particles to be scattered by
a thin metal foil. Thus, scattering process involved
simultaneously many particles, but we assume that
these particles among themselves exert no forces,
central or otherwise, while each of them experi-
ence a central force at the scattering centre. We
consider a beam of particles, whether electrons, α
particles or planets of the solar system is irrelevent
all of the same mass and energy incident upon a
centre of force. It will be assumed that the force
falls off to zero for very large distances.
The incident beam is described by specifying
its intensity I, also called flux density, which gives
the number of particles crossing unit area normal
to the beam in unit time.
Physical Sciences | 161
As the particles approaches the centre of force,
each will initially move in the beam direction, but
will deviate from it as it nears the centre. After
passing the centre of force, the particle will, at a
large distance, once again pass along a straight
line. In general the final direction of motion (away
from the centre) is not the same as the incident
direction and the particle is said to be scattered.
The cross section for scattering is a given
direction, σ (Ω) is defined by
Number of particles scattered into
solid angle d Ω per unit time
σ (Ω) d Ω =
Incident intensity
where d Ω is an element of solid angle in the
direction Ω.
The above called the differential cross section,
with central forces, there is a completely symmetry
around the axis of the incident beam. Hence, the
element of solid angle may be taken as
d Ω = 2π sin θ dθ …(1)
where θ is the angle between the scattered and
incident directions, called the scattering angle.
The name ‘cross section’ has been given
because the dimensions of σ (Ω) is that of an area.
We express the angular momentum of the particles
in terms of the energy and ‘Impact parameter’
s, defined as the perpendicular distance between
the centre of force and the incident direction of
velocity.
In v 0 be the incident speed of the particle,
then initial angular momentum
= l = mv0s = s
√
2mE …(2)
s θ
dθ
ds
Fig. : Scattering of a beam of particles by centre of
force
Once E and s are fixed, the angle of scattering
θ is then determined uniquely (In classical theory).
We assume that different values of s cannot lead
to the same scattering angle.
Hence, the number of particles scattered into
a solid angle d Ω lying between θ and θ + dθ,
must be equal to the number of incident particles
162 | Physical Sciences

with impact parameter lying between the Now consider the case of a repulsive Coulomb
corresponding s and s + ds. field. The change on incident particle is (Z′e) and
Hence, 2π ls | ds | = 2π σ (θ) l sin θ | d θ |. that at the scattering centre (target element) is Ze.
The number of particles must always remain The force is then given by
positive even if ds and d θ do not have the same ZZ′e2
sign. f = –
r2
Let s be considered as a function of energy E
The force constant
and the scattering angle θ, i.e., let s = s (E, θ) then,
a ds k = – ZZ′e2 …(8)
σ(θ) = …(3) The energy ε > 0 and the orbit is a hyperbola
sin θ dθ
with the eccentricity given by
We will for simplicity, consider the case of

√

repulsive scattering. An expression for the 2 El2
scattering angle θ as a function of s may be ε = 1+ 2 2
m (ZZe )
directly obtained from the orbit equation,
θ = ⌠

√ 1+(
ZZ′e )
2 Es
r 2
 dr
+ θ0 …(4) = …(9)
 2


√
 r2 2mE 2mV 1

⌡r0 l2
–
l
–
r2 Taking the periapsis direction to be θ = 0, the
orbit equation becomes
As the orbit must be symmetric about the 1 m (ZZ′e2)
direction of the incident beam, the scattering angle = (s cos θ – 1) …(10)
is given by r l2
θ = π – 2ψ …(5) The direction of the incoming asymptote ψ is
then determined by the condition r → ∞ which
where ψ is the angle between the incoming
gives
asymptote and the direction of incident beam
1
(Periapsis) cos ψ = …(11)
ε
By equation (5),
θ 1
sin = …(12)
2 ε
θ 2Es
We also obtain ψ from equation (4) by setting Hence cot2 =
2 ZZ′e2
r0 = ∞ when θ0 = π (the incoming direction) when ZZ′e2 θ
θ = π – ψ, r = rm the distance of closest approach This gives, s =
2E
cot2
2
…(13)
Then
From equation (3) we get
ψ = ⌠
∞
 dr 1 ZZ′e2 2 θ

…(6) σ(θ) = ( ) cosec2 …(14)
 r2 2mE

√
2mV 1 4 2E 2

⌡rm l2
–
l
– 2
r This is Rutherford scattering cross section
Expressing in terms of impact parameter s, formula for scattering of α particle by nuclei
under the central repulsive inverse square law
θ(s) = π – 2 
⌠r sdr force.
 r


√
We may define a total scattering cross

⌡r
( )
r2 1 –
V(r)
E
– s2 sections σT by

⌠ rm
m
σT = ∫ 4π
σ (Ω)d Ω
θ = π–2
rs lu π
∫
or …(7)
 σT = 2π σ (θ) sin θ d θ

√
or
 V(u) 2 2
 1– –s u 0
⌡0 s However, the calculation of σT for Coulomb
1 scattering by substituting (14) in this integral and
where we put = u. carrying out the integral, leads to an infinite value.
r
 Physical Sciences | 163

Coriolis Force Coriolis force acting on a particle moving

— — →
Coriolis force is given by – 2m (w × v a ) with a velocity v in a reference frame rotating
— → —
where v a is the velocity of the particle with with an angular velocity ω is given by F cor =
respect to rotating frame. — —
– 2m ( ω × v )
It is zero if the particle is at rest relative to the As the particle has been projected
rotating frame. horizontally, the vertical component of its velocity
Horzontal & vertical components of — ·^ ·^ ·^
coriolis force due to rotation of earth on a v, zk = 0 & horizontal components are x i + y j .
particle of mass m projected horizontally with To find the value of coriolis force evaluate
— —
a velcoity v ( ω × v ) which is given by
P → point on the surface of earth at lattitude
π
^i ^j ^k 
[
λ co-lattitude φ = ( )]
2
–λ ω × v = 
— —
0 ω y ω z 
·x ·y 0 


ω ω
y z = (– ω y) · ^i + (ω x)
· ^j – (ω ·x)k^
z z y
N x East → · sin λ ^i
P ∴ F cor = – 2m [– yω
λ + xω sin λ ^j – xω
· · cos λ k]
^
Equator
= – 2mω [–y· sin λ ^i + ·x sin λ^j
– ·x cos λ^k]
The magnitude of coriolis force is
S Polar axis →
| F | = 2mω [·y2 sin2 λ + x· 2 sin2 λ
cor
z-axis of the system is taken vertically
upward at P; radially outward from the centre of + ·x2 cos2 λ]1/2
the earth. = 2mω [y2 sin2 λ + x· 2]1/2
·
The XY plane is the horizontal plane = 2mω[vy2 sin2 λ + vx2]1/2
containing the point P, the x-axis pointing towards In terms of co-lattitude φ,
the east & y towards north. →
Earth rotates from west to east and the + ve | F cor | = 2mω [vy2 cos2 φ + vx2 ]1/2
direction of x-axis taken towards east. The The horizontal component of coriolis force
angular velocity vector lies in the XY plane in a is given by—
direction parallel to the polar axis about which →
→ ( F cor)horizontal
earth rotates. The vector ω , therefore, has no = –2mω[–y· sin λ^i + ·x sin λ^j ]
component in the East-West direction i.e., along
the x-axis = –2mω sin λ[–·y^i + ·x^j ]
The magnitude of horizontal component is
If ^i , ^j & ^k one the unit vectors along X, Y given by
and Z axes respectively, →
→ |( F ) | = 2mω sin λ[·y2 + x· 2]1/2
Then ω = ω x ^i + ω y ^j + ω z^k cor H
= 2mω sin λ[vx2 + vy2]1/2
where ωx = 0
= 2mωv sin λ
ω y = ω cos λ,
In terms of co-latitude φ, the magnitude of
ω z = ω sin λ horizontal component is
Let a particle be projected from P with a 2mω cos φ [vx2 + vy2]1/2
— = 2mv cos φ
velocity v ,
→ The horizontal component of coriolis force
Then v = v ^i + v ^j + v ^k
x y z has max value at the poles (λ = 90°) and zero at
= ·x^i + ·y ^j + ·z ^k the equator (λ = 0)
164 | Physical Sciences

⇒ The vertical Component of Coriolis Radius of gyration—If the entire mass of the
force is given by body is supposed to be concentrated at a point
→ · such that the K.E. of rotation is the same as that of
( F cor)vertical = –2mω[–x cos λ^k] body itself, then the distance of that point from the
·
= +2mωx cos λk^ axis of rotation is called radius of gyration of the
The magnitude of the vertical component is body about that axis.
given by 1
→ K.E. = Iω2
2
| F | = 2mω·x cos λ
cor v
1
= 2mωv x cos λ = Σmr2 ω2
2
In terms of co-latitude
= 2mωv x sin φ 1
= mk2ω2
→ 2
The vertical component ( F cor)vert = + 2m w ·x mk 2 = Σmr2
cos λk^ i.e., it acts in the direction of + z axis which r1 2 + r2 2 + …
means vertically upward. At the equator if the = mn { n }
body has an initial velocity x, · it will appear to be m = mn
lifted upward. r1 2 + r2 2 + … 1/2
Direction of Coriolis force in northern and
southern hemispheres
and k = ( n )
Theorem of perpendicular axis—It states
→ →
⇒ – 2m (ω × v ) ⇒ coriolis force i.e., that the moment of Inertia of a plane lamina about
→ an axis perpendicular to the plane of the lamina is
negative of the cross product of ω → angular equal to the sum of the moments of inertia of the
→
velocity of rotation of earth and v the linear lamina about the two axis at right angles to each
→ other, in its own plane intersecting each other at
velocity of the body. The direction of ω the
angular velocity vector is always along the axis of the point where the perpendicular axis passes
the earth from S to N. through it.
In the northern hemisphere if a body is I = Ix + Iy
moving towards north (i.e., along + Y-axis) an Y
application of the rule for the direction of vector
product of two vector indicates that this force will x1
act towards the east (i.e., along + X-axis). O n P
the other hand if the body is moving toward r1 x2
south (i.e., along -Y-axis) the coriolis force will O X
acts towards the west (i.e., along-X axis).
In other words, in the northern hemisphere a
moving body turns towards the right.
Theorem of parallel axes for moment of
Similarly in the southern hemisphere a moving
inertin—It states that the moment of inertia of
body will turn towards the left due to the effect of
a body about any axes is equal to the sum of
coriolis force.
its moment of inertia about parallel axis through
Moment of Inertia in centre of gravity and the product of its mass
The moment of Inertia of a body about an and the square of the distance between the two
axis is defined as the sum of the products of the axis h → perpendicular distance between two axes
mass and the square of the distance of the is h
different particles of the body from the axis of G
A
rotation.
1
⇒ K.E. of rotation = Iω2 m1
P
2 n
If ω = 1
the I = 2 × K.E. B
 Physical Sciences | 165

M.I. of circular disc about an axis through, MR 2 1

So, I1 = ×
its centre perpendicular to its plane— 2 2
M MR 2
Mass per unit area = =
πR2 4
(b) About tangent—M.I. of disc about dia-
R
meter
MR 2
=
x 4
dx
C
X

M R
Mass of element = × 2πxdx
πR2 A O
2M
= x dx
R2
M.I. = (M of element) about Y
B
an axis through its
centre ⊥ to plane Theorem of parallel axis
I = I1 + M × OA2
= ( 2M
R2
x dx x2) =
MR 2
+ MR2
2M 3 4
= x dx
R2 5
= MR2
About whole disc I = ∫ x dx
2M R 3
R2 0
4
Solid cylinder—
1 (a) axis of symmetry
= MR2
2 1
M. I. of circular disc : MR2
2
(a) about a diameter (b) about the axis passing through its centre
(b) about a tangent and perpendicular to its own axis of symmetry—
(a) About a diameter—M.I. of disc of a cir-
cular disc about an axis ⊥ to its plane and passing y1 l/2
through centre is
C O x D
y

A B
M.I. about any diameter
O M R2
= dx ×
l 4
M.I. about llel to YY′
D M R2 M
= dx + dx x2
1 l 4 l
I = MR2 M R2 2
2
⇒ Theorem of perpendicular axis
=
l 4 (
+ x dx )
I = I1 + I 2 l l
⇒ Complete it by integrating + to –
Two diameters AB, CD, symm w.r.t. disc 2 2
I1 = I2 M.I. of thin rod—
I = 2I1 ⇒ Hint I = Part b of solid cylinder
I and In this case R = 0
I1 = Also find K.E. and Ans. momentum.
2
166 | Physical Sciences

Table for M.I.

(1) Thin uniform Rod : (1) Through centre ⊥ to its length. Ml 2 /12
(2) Through one end and ⊥ to its length. Ml 2 /3
(2) This rectangular lamina Through the centre and ⊥ to its plane. M/12 (l2 + b2 )
(3) Hoop or circular ring (1) Through the centre and ⊥ to its plane. MR 2
(2) About diameter MR 2 /2
(4) Circular lamina or disc (1) Through the centre and ⊥ to its plane. MR 2 /2
(2) About dimaeter. MR 2 /4
(5) Angular ring or disc, of outer (1) Through the centre and ⊥ to its plane. M/2 (R1 2 + R22 )
and inner radii R2 and R1 . (2) About a diameter. M/4 (R1 2 + R22 )
(6) Solid cylinder (1) About its axis of cylindrical symm MR 2 /2
and about its own axis. 2 2

(2) Through its centre and ⊥ to axis. M (x1 + R4 )

2

(7) Hollow cylinder ext. radii R2 (1) About its own axis M
(R 2 + R22 )
int. radii R1 2 1
(2) Through its centre and ⊥ to its axis. R 1 2 + R2 2 l 2
M ( 4
+
12 )
(8) Solid sphere (1) About a diameter. 2/5 MR2
(2) About a tangent. 7/5 MR2
(9) Spherical shell (1) About a diameter 2/3 MR2
(2) About a tangent 5/3 MR2
(10) Solid right circular cone of (1) About its axis 3/10 MR2
base rad. R.
(11) Solid cone of altitude and About vertical axis 3/10 MR2
base rad. R
(12) Thick shell or hollow (i) About a diameter 2 R2 5 – R1 5 
5  R 2 3 – R 1 3 
sphere ext. radii R2 int. radii R 1 M

(ii) About a tangent 2  R 2 5 – R 1 5

5  R 2 3 – R 1 3
+ MR22

Mechanics of System of Particles Thus, the total force on the ith particle will be
A collection of large number of particles is ext int
F i = F i + Fi
called a system of particles whose motion can be
N
= F i + ∑ F ik
studied by generalising the result of one particle ext

motion. k=1
k≠i
Let us consider a system of N particles of N
masses m 1 , m 2 , …, mn respectively. Let the where Fi = ∑ Fik, Fik is the force of interaction
int
k=1
resultant forces acting on the various particles be k≠i
represented by F1, F2, F3, …, Fn respectively. Each acting on ith particle due to kth particle. Hence,
of these forces is comprised of two parts : the equation of motion of the i th particle may be
ext
(i) External forces Fi written as
N
∑ F ik
ext
(ii) Internal forces arising due to the interac- mi ˙˙
r i = Fi +
int k=1
tion of various particles on each other. viz. Fi k≠i
 Physical Sciences | 167

Now, we obtain the equation of motion for Torque

the whole system by summing up the above The total torque force system of particles
equation over-all the particles, about a point 0 is defined as
N
∑ mi ˙˙
N
r i = ∑ Fi + ∑ ∑ F ik
ext N N
→ → →ext
i=1 i=1 i=1 k=1 τ = ∑ ri× Fi
k≠i i

Since from Newton’s third law where ri is the position vector of the ith particle
from 0.
F i k = – Fki
Let R be the position vector of the centre of
Hence, the equation of motion for the whole mass of the system relative to 0 and r is the
system reduces to position vector of ith particle relative to the centre
N of mass as shown in figure below :
r i = Σ Fi = Fext
∑ mi ˙˙
ext
i=1 y
Now, if we introduce a point with position
vector R, such that
Σ miri 1 ith particle
R = = Σ miri
Σ mi M ri ri′
where Σ mi = M = total mass of the system. This
point with position vector is called the centre of Centre of mass
R
mass or centroid of the system of particles. Hence, x
O
the equation of motion for the whole system
reduces to Here ri = R + ri . Thus, vi = v + vi′
M˙Ṙ = F ext Hence, we have
τ = ∑ (R + r′i) × Fi = τi + τc
ext
This equation refers to the equation of motion
for a particle of mass M = Σ m i having position
i

where τ′i = ∑ r′ i × Fi
ext
vector R thus— is the total torque about
i
The centre of mass moves as if the total mass the centre of mass and
∑ R × Fi
ext
of the system were concentrated at it and the total τc =
external force were acting on it. i
= R × Fext
Linear Momentum d
= (P × MV)
dt
The total linear momentum of a system
particles is defined as is the torque of centre of mass about the origin O.
→ Angular Momentum
P = Σ mivi
It is defined as
= Σ mi ṙi = MṘ L = ∑i ri × pi
or
dp
= M˙Ṙ = Fext L = ∑i ri × mi ṙi
dt
Hence the rate of increase of linear momen-
= ∑i mi (ri × vi)
tum of a system of particles is equal to the vector N
sum of the external forces. Lx = ∑ mi (yi żi – ziẏi)
i=1
N
If Fext = 0 then P is a constant vector. Ly = ∑ mi (zi ẋi – xi żi)
i=1
Thus, if the net external force acting on a
N
particle vanishes, its linear momentum is con- Lz = ∑ mi (xi ẏi – yi ẋi))
served. i=1
168 | Physical Sciences

Other results are in this case as follows :  1 

 ∑i V i + 2 ∑i ∑k Vik 
ext
where V =
(i) L = L′ + Lc where L c = M(R × V) is the  
i≠k
angular momentum of total mass M located at the
centre of mass about O and is the total potential energy function for the whole
system.
L′ = ∑i mi (ri × v′i)
is the total orbital angular momentum of the
Conservation of Energy
system about the centre of mass.
We see that :
(ii) If the not torque acting externally on the
system is zero, the total angular momentum of the WAB = TB – TA = VA – VB
system remains conserved. or TA + VA = TB + VB = constant
(iii) The rate of change of angular momentum Showing that the sum of kinetic and potential
of a system about a point, either fixed or moving energies is conserved provided both the internal
with the centre of mass is equal to the total and the external forces are conservative.
moment of external forces about that point or in
symbols. Phase Velocity and Group Velocity
dL
= G The general equation for any plane progres-
dt
sive wave is given by
Kinetic Energy → →
y = a sin (ωt – k . r ) …(1)
Since energy is measured in terms of work
2π 2π
done, so we find the work done in displacing the where ω = is angular frequency and k = is
system from a position A to position B, equal to T λ
the difference between the kinetic energies at wave vector or propagation constant.
those points and is given by If the wave progress in + X direction
WAB = TB – TA y = a sin (ωt – kx) …(2)
= ∑i ∫ A
B ext
Fi dri + ∑ ∑
i k ∫
A
B F ik dri
= a sin k (ωk t – x)
i≠k

It can also be proved that the total kinetic 2π

= a sin (vt – x)
energy of a system of particles about O is equal to λ
the kinetic energy of the centre of mass assuming
total mass to be acting on it plus the kinetic In wave motion the particles in the medium
energy of motion about the centre of mass. vibrate about the mean position. The velocity a
moving particle is known as particle velocity and
Potential Energy dy
is given by . In wave motion not only the
dt
If the internal forces F ik between the ith and
particles moves but there is a disturbance also.
kth particles are central forces i.e., they act along This disturbance also move with a finite velocity.
the line joining the particles i and k and the e.g., if a crest or through is generated at one end of
external forces are conservative, then the work string it reaches the other end after a certain time.
done by the total force in displacing the system The velocity of disturbance is known as wave
from A to B is given by velocity or phase velocity.
B
∑i ∑k Vik 
1
WAB = – ∑ V i
ext Phase velocity is the ratio of angular
+
i 2
i≠k
A frequency to the wave vector.
ω
= [–V]B = VA – VB vp =
k
…(3)
A
 Physical Sciences | 169

The propagation velocity of a plane progres- This is an equivalent equation of equation (6)
sive wave is the velocity of plane of constant whose amplitude is modulated by function
phase (wave front), that is why it is called phase
velocity. 2a cos (∆ω2 t – ∆k2 x)
For a wave front the phase is constant.
ωt – kx = constant …(4)
or ωdt – kdx = 0
dx ω
or = = vp …(5)
dt k
The velocity by which a single wave pro-
pagates is called as phase velocity. Whereas “the
velocity of wave packet formed by a group of
harmonic wave is known as group velocity” which
is equal to ( )
dω
dk
.

Let there is a group of wave having large Fig. : Modulation of two waves having a slight
difference in their frequencies.
number of waves having frequencies between ω
and ω + ∆ω or wave vector k and k + ∆k. If these Equation (9) represents a wave of maximum
waves superimposes, at some points of same amplitude 2a, frequency ω and wave vector k
phase the amplitude will be maximum whereas ∆ω
superimposed by frequency and wave
opposite phase will vanish the amplitude. There 2
will be points of reinforcement. These points of ∆k
reinforcement will constitute a finite region or vector , or this modulated wave is covered by
2
wave packet. The velocity of this packet is group ∆ω ∆k
velocity. an envelope of frequency and wave vector .
2 2
For convenience we take two waves having This envelope is nothing but a series of wave
amplitudes a but differs in frequencies by ∆ω and packets whose velocity is
wave vector by ∆k, i.e., ∆ω
vg = …(10)
y1 = a sin (ωt – kx) …(6) ∆k
y2 = a sin {(ω + ∆ω)t – (k + ∆k)x} …(7) For small variations is frequency and wave
vector
After super position
∆ω dω
y = y1 + y2 lim∆k → 0 =
∆k dk
1
= 2a sin [(2ω + ∆ω)t – (2k + ∆k)x ] or vg =
dω
…(11)
2 dk
1
cos (∆ωt – ∆kx) …(8) The component waves of this envelope may
2 have higher or lesser velocities. But the wave
∴ ∆ω << ω packet moves with group velocity.
and ∆k << k
so 2ω + ∆ω ≈ 2ω
2k + ∆k ≈ 2k
Hence, equation (8) be

y = 2a cos (∆ω2 t – ∆k2 x) sin (ωt – kx) …(9) Fig. : Wave packet
170 | Physical Sciences

Now, ω = k vp Hence phase may move with velocity higher

d than c but the group cannot.
∴ vg = (ω)
dk In easiest form ω = ck
dvp vg = vp = c
= vp + k
dk
Case III. vg > vp
dvp
or vg = vp – λ …(12)
dλ
or wave velocity not only depends on phase
(when dvdλ is negative)
p

velocity but also on the variation of phase velocity This is known as anomalous dispersion.
with wavelength.
vg > vp (Anomalous dispersion)
dispession)
dvp
Case I. = 0
dλ vg =vp (No dispersion)
dispession)
then vg = vp

ω (k)
This is the case of no dispersion. Hence, we vg < vp (Normal
dispession)
dispersion)
can say that in non-dispersive media the waves of
different frequencies moves with same group and
phase velocities. k

Case II. vg < vp ( dv

generally p is positive.
dλ ) This is to be noted that for dispersive media
vp is not significant physically. For non-dispersive
This is the case of normal dispersion e.g., media.
propagation of electromagnetic waves in refrac- ω ∝ k
ω ω = kv
tive media. For the media ≠ constant
k
v being the velocity of propagation.
vg ≠ vp
For dispersive media vg is only important and
take the example of propagation of X-rays in the transfer of energy takes place with this
carbon of refractive index n. velocity.
A Speed in vacuum
n = 1– 2
ω Also in n =
Speed in medium
,
ω A
Wave vector k = – The speed of wave is considered as vp.
c ωc
v v
ω
( Qk=
e/n
The velocity of light in
n = =
v′ vp
v
e
medium = not e‚ also λ =
n
e
n ) or vp =
dvp
n
v dn
Here A is a constant. ∴ = – 2
dλ n dλ
ω
Obviously vp = >c 1 dn
k = – vp
n dλ
dk 1 A
Now,
dω
= + 2
c ω c ∴ From equation (12)

=
1
c ( A
1+ 2
ω ) (–v ndλ
vg = vp – λ
dn
) p

λ dn
= v (1 +
n dλ)
dω c
∴ vg = = <c p
dk A
1+
ω2 where λ is wavelength in vacuum.
 Physical Sciences | 171

Special Theory of Relativity In the absence of any indication to the

contrary in our everyday experience, we further
Whenever we speak of ‘motion’ of course, we assume that
really mean “motion relative to a frame of
reference”. The frame of reference may be a road, t′ = t
earth’s surface, the sun, the centre of our galaxy; These set of equations are known as Galilean
but in every case we must specify it. There is transformations.
nothing to be called in absolute rest or absolute Although Galilean transformations are in
motion. accord with our intuitive expectations, these
The special theory is based upon two violates both postulates of special relativity e.g., if
postulates— we measure the speed of light is x direction in the
system to be c, however, is the S′ system it will be
(i) The laws of physics may be expressed in
equations having the same form in all frames of dx′
= vx – v or c′ = c – v
reference moving at constant velocity with dt
respect to each other. which is a violation of second postulate.
This postulate expresses the absence of an Lorentz Transformations
universal frame of reference. To overcome the disparity appeared between
(ii) The second postulate of special theory relativity postulates and Galilean transformations
of relativity states that the speed of light in free Lorentz developed a set of transformations directly
space has the same value for all observers, from the special relativity.
regardless of their state of motion. x′ = k (x – vt) …(1)
The Galilean transformations—Let us k is a function of v but independent of x and t.
suppose that we are in a frame of reference S and
find that the co-ordinates of some event that The inverse transformation be
occurs at the time t are x, y, z. An observer located x = k (x′ + vt′) …(2)
in a different frame of reference S′ which is The factor k will be same in both the frames S
moving unaccelerated with a velocity v with and S′, since there is no difference between S and
respect to S notice the event at time t′ and position S′ other than the sign of v.
x′, y′, z′. Let us say S′ is moving in +x
direction. The measurements x, y, z, t related to x′, y S y'
y′, z′ t′ are called Galilean transformations. P S'
r (x, y, z, t)
y y' r' (x', y', z', t')
S S' O x O' x'

O x O' x'
v z z'
νt
z z' There is no difference in co-ordinates along y
If time in both system is measured from the and z directions, i.e.
instant when the O coincides with O′, measurement y′ = y …(3)
in the x direction made in S will exceed those z′ = z …(4)
made in S′ by vt, which represents the distance The time co-ordinates are, however, not
that S′ has moved in the x-direction, that is equal. We can substitute the value of x′ from
x′ = x – vt equation (1) into equation (2)
There is no relative motion in y and z direction x = k2 (x – vt) + kvt′
So From which we find
y′ = y 1 – k2
z′ = z
t′ = kt + ( )
kv
x …(5)
172 | Physical Sciences

These equations constitute a co-ordinate trans- and inverse transformations be


formation that satisfies the first postulate of
relativity. x′ + vt
x =

 
√
The second postulate enables us to evaluate v2
k. At t = 0, O coincides with O′. The observers in 1– 2
c
both frames proceed to measure the speed with


which the light from it spreads out. Both will find y = y′ Inverse Lorentz


the same speed c. z = z′ transformations


In S, x = ct …(6) vx
while in S′, x′ = ct′ …(7) t′ +

1 – vc 
c 2
t =

√
Substituting for x′ and t′ in (7) with (1) and (5) 2
1 – k2
k (x – vt) = ckt + ( ) kv
cx
2

when the relative velocity v of S and S′ is small

and solving for x compared to the speed of light, Lorentz transfor-
ckt + vkt mations reduced to Galilean one.
x =
1 – k2
k– ( )
kv
c Addition of Velocities
 k + vk  Let us consider something moving related to
  both S and S′. An observer in S measures its
= ct  
c
 velocity components to be
k – c
2

 ( )1 –
kv
k
 Vx =
dx
dt
 v 
 1+  dy
= ct  c  Vy =
 dt
 1 – 12 – 1 c 
( )
 k v and Vz =
dz
dt
This equation is same as equation (6), namely while to an observer in S′ they are
x = ct so,
dx′
v V′ x =
1+ dt
c
= 1 dy′
V′ y =
( )
1– 2–1
k
1 c
v
dt
dz′
1 V′ z =
and k = …(8) dt


√
v2 By differentiating the Lorentz transformation
1– 2 equations, we obtain
c
So the set of transformations will be dx – vdt
dx′ =

 
√
v2
x – vt 1– 2
x′ = c

 
√ v2
1– 2
c
dy′ = dy
dz′ = dz
 v

y′ = y Lorentz dt – 2 dx
c

z′ = z transformations dt′ =


√
v2

vx 1– 2
t– c

1 – vc 
c 2
t′ = dx′ dx – vdt
Vx′ =
√
2 and So =
dt vdx
2 dt – 2
c
 Physical Sciences | 173

dx Squaring both sides, we get,

–v
dt V –v m 20
Vx′ = = x …(1) m2 =
vV x V v2
1– 2 1 – v 2x 1– 2
c c c
dy′ or m2c2 = m 20 c2 + m2v2
Again Vy′ =
dt′ Differentiating,

√ √
v2 v2 dy c2 2m dm = m2 2v dv + v2 2m dm
dy 1– 1–
c2 c2 dt or c2dm = mvdv + v2 dm …(5)
= =
v v dx From equations (3) and (5),
dt – 2dx 1– 2
c c dt dEk = c2dm …(6)

√vV
v2 Thus a change in K.E., dEk is directly propor-
Vy 1– tional to a change in mass dm. When a body is at
c2
= …(2) rest, its velocity and hence K.E. = 0 and m = m 0.
x
1– When its velocity is v, its mass becomes m.
c2
Therefore, integrating equation (6), we obtain

√v1 – vc
2

∫
Ek
Vz 2 Ek = dEk
Similarly Vz′ = …(3) 0

∫
m
1– Vx or Ek = c2 dm
c2 m0
Equations (1), (2) and (3) constitute the = c2 (m – m0)
relativistic velocities transformations. ∴ Ek = mc 2 – m0c2
Mass Energy Equivalence This is relativistic formula for kinetic
energy. When the body is at rest, the internal
We know that the force is defined as the rate energy stored in the body is m0c2 . m 0 c2 is called
of change of momentum i.e. the rest mass energy associated with the body.
dP The total energy (E) of the body is the sum of
F =
dt K. E. (Ek) and rest mass energy (m0c2)
or F =
d
(mv) …(1) ∴ E = Ek + m0c2
dt = (mc2 – m0c2) + m0c2 = mc2
According to the theory or relativity, both or E = mc 2
mass and velocity are variable. Therefore, This is known as Einstein’s mass energy
d equivalence formula. Which means that mass and
F = (mv)
dt energy are the two aspects of same thing.
dv dm The formula for K. E. reduces to the classical
or F = m +v …(2)
dt dt formula for v << c
Let the force F displace the body through a Ek = mc 2 – m0c2
distance dx. Then, the increase in kinetic energy = (m – m0)c2
(dEk) of the body is equal to the work done F dx. v2 –1/2
Hence, dEk = Fdx
= m0c2 [( )
1– 2
c
–1 ]
dv dm Now since v << c
= m dx + v dx v2 –1/2 v2
or
dt
dEk = mvdv + v2dm
dt
…(3)
( )1– 2
c
= 1 + 2 + ……
2c
(using Binomial expension and leaving higher
According to the law of variation of mass
order terms.)
with velocity,
v2
m0 ∴ Ek = m0c2 × 2
m = …(4) 2c


√
v2 1
1– 2 = m0 v 2
c 2
 Unit-3

ELECTROMAGNETICS
Electric Charge
Electric charge is a scalar quantity which
remains conserved in space and time. The quanta
of electric charge is ‘e’, 1·6 × 10–19 coulomb. The
net charge of an isolated system remains
conserved.
where ^r ij is the unit vector along the direction of
line joining the two charges. The force is
attractive or repulsive depending upon the polarity
of charges. And ε0 is the permittivity of free
space, i.e.,
ε0 = 8·85 × 1012 C/N-m2

Charge Densities Electric Field Intensity

For the continuous distribution of charge in The region surrounding any point charge in
space, we can define the charge densities as which the presence of this charge can be
follows : experienced is known as electric field. The
intensity of this field at any point is defined by a
(a) If ∆q is the net charge enclosed by a vector function E, which equals the force
volume ∆V the volume charge density be experienced by unit positive charge
∆q →
ρ = lim∆V → 0 →
∆V F
E = limq0 → 0
(b) If the charge ∆q is spread over a surface q0
are ∆S the surface charge density be → 1 q ^
or E = .r
∆q 4πε0 r2
σ = lim∆S → 0
∆S The unit of electric field intensity is N/C or
(c) If the charge ∆q is distributed along a line volt/metre Vm–1. And it is some times referred as
then linear charge density be field strength.
∆q Electric Potential
λ = lim∆l → 0 .
∆l The electric field intensity due to a point
The total charge in a finite region can be charge q is
given by →
→
q = ∫ dq = ∫V ρ dV + ∫S σ dS E =
q r
4πε0 r3
Coulomb’s Law
Coulomb discovered that the electrostatic
Also ∇ ()
1
r
= – 3 ()
r
r
→
force between two point charge varies in inverse
square manner and acts along the line joining two
Thus E =
q
4πε0 [ ( )]
–∇
1
r
stationary charges.
The force between the charges qi and qj will
= –∇ ( ) q 1
4πε0 r
→
be or E = – ∇V
– 1 qi qj ^ 1 q
F ij = r V = .
4πε0 r2ij ij 4πε0 r
224 | Physical Sciences

The electric field can be described in terms of using r << r0 and expanding in binomial fashion,
a scalar function known as electric potential.
Electric potential is a scalar quantity whose
negative gradient provides us the electric field.
[ 1 3
(1 + x)1/2 = 1 – x + x2 – …
2 8 ]
1 ⌠  ρ  1 + r cos θ + r22
The unit of electric potential is volt. φ =
4πε0 r0 
⌡  r0 r0

(3 cos2 θ – 1) + … ]dV
2

4πε r ∫
1 1
= ρ dV + 2
0 0 4πε r 0 0

∫ ρ r cos θ dV + 4πε1 r 0 0
3

∫ ρ r (3 cos2 θ – 1)dV + …
2
2

K0 K1 K2
= + + +…
4πε0 r0 4πε0 r 20 4πε0 r 30
Fig. : Figure showing the relations between potential
V electric field E and charge density ρ. Here K0 = ∫ ρ dV
The Moments of a Charge Distribution K1 = ∫ ρ r cos θ dV

∫ 21 ρr (3 cos2 θ – 1)dV
2
Let there is an arbitrary charge distribution as K2 = 2
shown in the figure. Due to which we have to
calculate the electric field at a point P far from
charge distribution.

o Hence potential at ρ can be expressed in the

ro
r 1
power series of . The value of Ks depends on the
r0
R
symmetry of charge distribution. K0 = ∫ ρ dV is
A P
dV
The potential at point P due to an infinitely known as monopole moment. For equal positive
small elementary area dV. and negative charge K = 0.
1 ρ dV The second term K1 = ∫ ρ r cos θ dV represent
dφ =
4πε0 R the dipole moment term. In which small positive
here ρ dV is the charge confined in small volume and negative charges are separated by small
element dV. distances. This term is effective only when K0 = 0.
The third term is a representation of
∴
φ =
1
4πε0 ∫ ρRdV quadrupole moment. This term is effective when
K0 and K1 are zero. The potential due to
1
In ∆ OAP quadrupole moment is proportional to 4.
r
R = (r20 + r2 – 2rr0 cos θ)1/2 Next term will represent octopole and so on.
(from cosine law) – +

⌠
+ – + –

φ =
1  2 2 ρ dV
4πε0 
⌡ (r0 + r + 2r0 r cos θ)1/2
+ –

+ – + – + –
+

∫ ρ[1 + (rr ) – 2(rr ) cos θ]

1/2 Monopole Dipole Quadrupole Octopole
1 2
= dV 1
( V ∝ —)
1
( V ∝ —)
1
( V ∝ —)
1
( V ∝ —)
4πε0 r0 0 0 r r2 r3 r4

Gauss’ Law
This is a law of prime importance in electro-
statics which relates the electric flux with charge
enclosed.
The electric flux φ E through any closed
1 supposed to be concentrated at a point, say the
surface is equal to times the net charge q
ε0
enclosed by the surface, i.e.
q supposed to act as a point, say centre of gravity of
φE = .
ε0
The law can also be expressed as integral of
the normal components of electric field (directed
outward) over a closed surface of any shape in
1 are displaced from each other, the molecule as a
electric field times the net charge enclosed by
ε0
the surface, i.e.,
S
∫
O = E . n^ dS =
q
ε 0
∫ q
or O = E dS =
S
ε0
∫
Here O represents the surface integral over
S Polarization
whole of the closed surface and n unit vector
drawn normal to the surface element dS.
Dielectric Polarization
The materials whose electrons are not free to
move about under the influence of an electric field
are known as insulators. The presence of the
electric field may change the behaviour of an
insulator. If the constitution of insulator gets
modified in the presence of electric field are
called dielectric. When the change in behaviour of
dielectric is independent of direction of electric
field applied it, is called isotropic on the other
hand. If the change in behaviour of dielectric
depends on the direction of field applied it is
called anisotropic.
The electric field exerts a force on every
charged particle of dielectric. The positive and
negative charges are pushed in opposite
directions. As a result the positive and negative
sectors of a molecule shifts from their equilibrium
positions. This shifting is limited by restoring
forces set up by the changing charge configuration
of the atom or molecules.
This relative displacement or shifting of
centre of charges is termed as polarisation and the
dielectric is said to be polarised.
Physical Sciences | 225
On the basis of polarization the dielectrics are
categorized in two classes :
(i) polar molecules and
(ii) non-polar molecules.
The positive charge of the nucleus may be
centre of gravity of positive charge. Similarly, the
negative charge due to orbital electrons may be
negative charge.
When the two centres of gravity coincide in a
molecule the molecule as a whole possesses no
resultant charge and is said to be non-polar, e.g.,
carbon tetrachloride. If the two centres of gravity
whole possesses polarity and has permanent
electric moment and the molecule is said to be a
polar molecule, e.g., NH3 , C 2 H5OH etc.
A polar molecule in an electric field exhibits
both permanent and induced dipole moments
whereas non-polar molecules exhibits
induced dipole moment only.
In a homogeneous isotropic dielectric the
polarisation (Reorientation of electronic orbits)
depends directly on the total electric field E in the
dielectric medium.
For small value of total electric field E. The
average dipole moment P induced in each
molecule of the dielectric is proportional to
electric field E.
P ∝ E
P = α E,
here the proportionality constant α is known as
molecular polarizability. For n molecules per unit
volume
P = nαE
here ρ is electric moment per unit volume called
polarization.
A rectangular block of dielectric is shown in
the figure whose length is in the direction of
electric field. The surface charge density on the
surface of dielectric material is σp.
The induced dipole moment be σp . a. l,
where σp . a is charge.
226 | Physical Sciences

By definition total induced dipole moment From Gauss’ law

= P.a.l ρ
⇒ P . a . l = σp . a. l ∇.E = …(1)
ε0
or P = σp here ρ is the charge density of region.
Polarization can be defined as the surface In a dielectric of infinite extent,
charge density appearing at the faces perpen-
dicular to the direction of applied field. ρ + ρ′
∇.E = …(2)
Again, P ∝ E For small fields ε0
or P = χe ε0 E here ρ and ρ′ are free change density and bound
here χ e is a dimensionless quantity known as charge density respectively.
electric susceptibility. Hence ∇ . ε0 E = ρ + ρ′
Also ρ′ = – ∇ . ρ …(3)
Dielectric Constant and Electric Dis-
placement Vector D So (3) became
∇ . ε0 E = ρ – ∇ . ρ
The capacitance of a parallel plate capacitor
increases when a dielectric substance is filled or ∇ . (ε0 E + ρ) = ρ
between its plates. ∇.D = ρ
The dielectric constant (K) or relative This is called Gauss’ dielectric law here
permittivity (ε r) is the ratio of capacitance of
D = ε0 E + ρ.
parallel plate capacitor filled with dielectric to the
capacitance without dielectric.
Electrostatic Energy
ε Absolute permittivity of dielectric
εr = = Consider the electric field arising due to some
ε0 Absolute permittivity of free space
distribution of charge.
Also we can show that
The work done in assembling the charge from
εr = 1 + χe infinity to establish the given distribution is called
q 2 the electrostatic energy of the field. The energy
r
4 πε0 ε E density is expressed as
Again, εr = = ⇒ 0
q 2 ε0 E 1
r U = E.D
4πε 2
The electric field intensity depends upon the 1
medium, we define another electrical vector or U = ε E2
2
which is independent of nature of medium.
1 q → Boundary Value Problems in Electro-
D = r statics
4π r3
 q → r 
For simple charge distributions the electric
D = ε  →
4πε r3 
or
field vector E and potential V at arbitrary point in
space can be found either by direct integration
⇒ D = εE from Coulomb’s law or by Gauss’ law. The
c boundary value problems in electrostatics consist
D is analogous to E having dimension . It
m3 of finding potential V for complex fields which
is called electric displacement vector. The product satisfies the specified boundary conditions.
D . ds gives the flux of normal displacement The separation of variables, image method
through surface ds and remains unchanged during and Green function etc. can be used to find V.
medium transfer. In an isotropic homogeneous Once V is known electric field, charge distribu-
field the directions of both vectors coincide. tion and total charge on conductors etc. can be
Relation between E and D. deduced.
 Physical Sciences | 227

Boundary Conditions
The solution of boundary value problems
should satisfy the following conditions at
boundary :
(i) Electric potential V must be continuous
across any boundary. Because non-existence of its
derivative means an infinite field, which is
physically not possible. Also the potential at
infinity must tends to zero for finite distribution of
charge and it must be constant throughout the
conductor.
(ii) Normal component of electric dis-
real nor can be imaginary. When two are real the
third will be imaginary and vice-versa.
The imaginary constants lead to oscillating
solutions whereas real constants lead to expo-
nential solution.
Method of Separation of Variables in
Spherical Polar Co-ordinates
Laplace equation is spherical polar co-
ordinates can be given by
1 ∂
∇2 V =
r2 ∂r(r ∂V∂r ) + r sin1 θ ∂θ∂
2
2

(sin θ ∂V∂θ ) + r sin1 θ ∂∂φV = 0

placement vector D should be continuous across a 2
charge field boundary. D1 n = D2n . And if the 2 2 2
boundary is between a conductor and a dielectric
say medium 1 is dielectric and medium 2 is to find the solution, we assume
conductor then D1n = σ and D2 = 0. V (r, θ, φ) = R(r) Θ(θ) Φ(φ ).
(iii) For boundary between two media the Substituting in above and dividing by RΘΦ
tangential component of electric field intensity we get
must be continuous across the boundary, i.e.,
sin 2 θ ∂ ∂R sin θ ∂ ∂Θ
E1t = E2t R ∂r
r2 ( )
∂r
+
Θ ∂θ (
sin θ
∂θ )
and if the boundary is between a conductor and a 1 ∂2 Φ
dielectric E1t = E2t = 0. So electric field intensity = –
φ ∂ φ2
is always normal to the surface of a conductor.
L.H.S. of this equation is a function of r and
Methods to Solve Boundary Value θ whereas R.H.S. is a function of φ only.
Problems ∂2 Φ
Let = – m2 φ having solution
(a) Method of separation of variables in carte- ∂ φ2
sian co-ordinates. The Laplace equation ∇2V = 0 Φ = C me± imφ
in cartesian co-ordinates is given by
and
∂2 V ∂2 V ∂2 V
1 ∂
(r ∂R∂r ) = – θ sin1 θ ∂Φ∂ (sin θ ∂Θ
2

∂θ ) sin θ
+ + 2 = 0 m
∂x2 ∂y2 ∂z 2 +
R ∂r 2
to find a general solution of this equation by the
method of separation of variables we assume that The variables are separated if we equate each
side of equation to constant n (n + 1) giving
V(x, y, z) = X(x) + Y(y) + Z(z)
∂ ∂R
when this is substituted in above Laplace equation
we get equations in which every term is
∂r ( )
r2
∂r
– n (n + 1)R = 0

independent of other two and can be given as and

1 ∂ ∂Θ m2
1 ∂2 X
X ∂x2
= c1 sin θ ∂θ (
sin θ
∂θ ) [
+ n (n + 1) – 2
sin θ ]=0

1 ∂2 Y and yield solutions

= c2
Y ∂y2 B
R = A n rn + n +n 1
1 ∂2 Z r
= c3
Z ∂z2 and Θ = P n m (x)
with c1 + c2 + c3 = 0. = Associated Legendre function
It is obvious from this auxillary condition that 1 dn 2
where Pn (x) = n (x – 1)n and x = cos
the separation constants c1, c 2 and c3 cannot all be 2 n! dxn
228 | Physical Sciences

Thus, the complete solution of Laplace Biot-Savart law is

equation be
⌠→ →
µ0 i  dl × r
v (r, θ,φ) = [A n rn + Bn r–(n + 1)] Pn m (cos θ) 
4π ⌡ r3
B =
C m e± imφ
→
where A and B are arbitrary constants which can
be determined by boundary conditions.
and
r
r3
= –∇
1
r
, ()
Biot-Savart Law where ∇ is three dimensional differential operator
∂ ^ ∂ ^ ∂ ^
This law is magnetic equivalent of Coulomb’s ∇ ≡ i+ j+ k
law. Biot-Savart studied the magnetic field ∂x ∂y ∂z
generated near a current carrying conductor and So the Biot-Savart law became
various reasons to establish the law.
According to this law the magnetic field
generated by an elementry length ∆l of a current
B = – ∫ dl × ∇ (1r )
µ0 i
4π
carrying conductor,

4π ∫ (r )
µ i 1
i ∆l sin θ = ∇
0
× dl
∆B∝
r2
where r is the distance of point of reference from
the elementry length [(→A ×→B ) = (→B ×→A )]
Again
i
∇ (1r ) × dl = ∇ × (dlr) – r1 (∇ × dl)
∆l

θ
r P
(By vector identity)
So,

∫ (∇ × dlr)– ∫ r1 (∇ × dl)
µ0 i 
In vector form it will be 
4π  
B=
⌠→ →
µ0  dl × r 
i
4π ⌡ r3
B =
dl is independent of (x, y, z), i.e., ∇ × dl = 0
µ0
µ0 i
∫ r1 ∇ × dl
where is constant of proportionality.
4π B =
4π
The major difference between Coulomb’s law
in electrostatics and Biot-Savart law in magnetism or by interchange of operation
is the direction of field. In Biot-Savart law the
→
direction of induced field d B is perpendicular to
the plane containing d l and r whereas in
B =
µ0 i
4π ∫ dlr
∇×

B = ∇× ∫
Coulomb’s law the electric field direction is
 µ i dl 

0
always in the plane containing charges.
 4π r 
Magnetic Vector Potential
B = ∇×A
In electrostatics we are familiar with relation
E = ∇φ, or the electric field is negative gradient of
potential in similar fashion magnetic field B is
where A =
µ0 i
4π ∫ dlr.
related to a quantity A by the relation B = ∇ × A,
where A by analogy is called magnetic vector The magnetic induction B is given by the curl
potential. of magnetic vector potential A.
 Physical Sciences | 229

Magnetic Scalar Potential Double Refraction

If we choose a scalar potential such that Some crystal do not obey Snell’s law of
B = – ∇ . φ m , φ m is said to be magnetic scalar refraction, in these crystals—
potential (i) There are two refracted rays for each
µ incident ray
φm = 0 Iω
4π (ii) Two refracted rays may not lie in the
here I is the current flowing through given plane with incident ray.
reference coil and ω is the solid angle subtended (iii) Each of the refracted ray is plane
by the coil on point of observation.
polarised. The polarisation in one being perpen-
Maxwell’s Equations dicular to that of other.
Maxwell unified the electric and magnetic This phenomenon is known as double refrac-
phenomenon with the help of four equations these tion and the crystals are known to be optically
are anisotropic. In some crystals one refracted ray
ρ follow the refraction law (ordinary ray) whereas
∇.E = (Gauss’ law) other violets the law of refraction (extraordinary
ε0
ray). If both the refracted rays violets the law of
∇.B = 0 (Absence of an isolated refraction the crystal is known to be biaxial
magnetic pole) crystal.
∂B
∇×E = – (Faraday’s law) Boundary Conditions at the Interface
∂t
∇ × B = µ0 J (Ampere’s law) between Two Media for e/m Field
Maxwell found that the Ampere’s law is Vectors
incomplete. In the space between the plates of a When any e/m wave incident on any interface
parallel plate capacitor in an A.C. circuit there is between two media then its vectors B, H, E and D
no conventional current (J = 0) yet ∇ × B ≠ 0. He, follow few boundary conditions—
∂∇ (i) The normal component of magnetic
therefore, introduced the term in addition to
∂t induction is continuous on interface, i.e.
conventional current term in fourth equation. B 1n = B 2n.
Now the modified Maxwell’s equation be
(ii) On the interface the tangential component
∇.D = ρ of electric field also remains continuous
∇.B = 0 E1t = E2t.
∂B (iii) The tangential component of magnetic
∇×E = –
∂t field also remains continuous
∂D H1t = H2t.
and ∇×H = J+
∂t (iv) The normal component of electric
Where J = ρ E, B = µ H and D = ε E displacement vector does not remain continuous
but their difference D1n – D2n always equals the
Maxwell’s equations can also be written in
surface charge density of interface
their integral form
D1n – D2n = σ.
∫
O D . ds = q
Reflection and Refraction of Plane EM
∫ B . ds
O = 0 Wave at a Plane Boundary of a
Dielectric
∫
O (∇ × E)ds = – ∫ ∂B
∂t
. ds When any e/m wave say light falls on any
interface a little fraction of it is reflected and
∫
O B . dl = ∫ (J + ∂D
∂t )
ds remaining fraction got refracted. Reflected and
refracted e/m waves follows the laws of reflection
Note : In free space ρ = 0 and εr = µr = 1. and refraction.
230 | Physical Sciences

Let any plane e /m wave incident on the or Kix . x + Kiy . y = Krx . x + Kry . y
interface at an angle θi as shown in the figure. The = Ktx . x + Kty . y
coefficients of refraction (Refractive index) of or Kix = Krx = Ktx
media are n1 and n 2 and K i, Kr and K t are the
and Kiy = Kry = Kty
propogation vectors of incident, reflected and
transmitted waves. Let the vector is in XZ-plane in XZ plane Kiy = 0
then the magnetic vector will be perpendicular to Hence reflected and transmitted waves are in
this plane. Vectors E i, Er and E t changes plane of incident wave.
periodically, so we can write Again
→ → →→ Kix = Krx = Ktx
Ei = Ei0e i (Ki r – ωit ) Ki sin θi = Kr sin θr = Kt sin θt
→ → → →
Ki and K r are same medium hence
Er = Er e i (Kr r – ωrt )
0 Ki = Kr
→ → →→
Et = Et e i (Kt r – ωtt ) or sin θ1 = sin θr
0
where Ei0 , Er0 and Et0 are the amplitudes of θi = θr
incident, reflected and transmitted waves and ωi, That is the law of refraction.
ωr and ωt is their frequencies. Dynamics of Charged Particles in
Static and Uniform Electromagnetic
Fields :
● Charged Particles Moving in Electric
Fields
If a charged particle is set free in an electric
field, it is accelerated by a force proportional to
electric field and charge on the particle as given
by
— —
F = e F (N) …(1)
According to Newton’s Second Law, the rate
of change of momentum (mv) of the particle equal
—
According to boundary condition the applied force F or
(Ei)tangential + (Er)tangential = (Et)tangential — d dv dm
F= (mv) = m + v (N) Newton’s
(Ei)x + (Er)x = (Et)x dt dt dt
Putting the value of Ei, Er and E t Second Law)
→→ → → …(2)
i (Ki r – ωit )
(Ei0)xe + (Er0)x e i (Kr
r –ωr t )
where m = mass of particle, Ks, v → velocity of a
→→ particle.
t – ω tt )
= (E t0)x e i (Kt If the velocity is very small compared to the
This holds good according to boundary velocity of light (v << c), the rate of charge of
condition. If the coefficients of r, y and t are mass is small (mass essentially constant), the last
equal, i.e. rerm in (2) is zero.
—
ωi = ωr = ωt = ω (let) So
—
F = m
dv
= ma
(coefficient of t) dt
So ma = eE
Hence, the frequency remains invariant dur-
ing reflection and transmission. The energy W acquired by the particles
∫
2 –
→ → → → → → W = m 1 a·dL
Ki . r = Kr . r = Kt . r
∫
2 – –
(coefficient of x and y) = e 1 E·dL = eV
 Physical Sciences | 231

V = ∫ E·d
– –
L
● The path of an electron in the transverse
deflecting fields is a parabola.
v =
dL ● The electron is subjected to the deflecting
dt
⇒
–
dL = vdt
field Ed = ( )
Vd
d
for a distance L or for a
L
dv time t = .
a = vx
dt
The field E d produces an acceleration a y in
v2 → final velocity the y-direction which is
v1 → initial velocity may = eE
∫
2
W = m V·dv eV d
1 =
d
= eV eV d
1 ⇒ ay =
– m(v22 – v12 ) = eV md
2
vy = ayt
If the particle starts from rest, the initial eV d L
velocity is zero, so =
md vx
1
W = eV = mv2 ⇒ particle energy The deflection angle α is then
2 v
α = tan–1 y
and particle velocity vx
eV d L
2eV (ms–1) particle velocity = tan–1 [1/2mvx2 = eVa]
v= mv x2 d
m VL
or, α = tan–1 d …(i)
For the electron 2Vad
| e | = 1·6 × 10–19 C From the tube geometry, assuming x >> L,
m = 9·1 × 10–31 C the angle α is also given by
y
v = 5·9 × 105 
√V α = tan–1 …(ii)
x
So, the deflection distance of the electron
The (CRT) Cathode Ray Tube, Electri- beam at the screen as
cal Deflection
VL
y = d x (m) Deflection distance
In CRT Cathode Ray Tube, charged particles 2Vad
(electrons) are accelerated and deflected so as to
strike a fluorescent screen in a desired pattern and Vd 2Vad
⇒ =
produce a visible image, y Lx
● A positive potential Va is applied to the Knowing this ratio, we can determine, what
accelerating electrode, or electrongen, it voltage is required to deflect the beam a
maintains this velocity v x in the vacuum of desired distance.
the tube, Charged Particles Moving in a Static
vy v Magnetic Field
Accelereting vd If a positively charged particle is set free in
electrode vx y —
Ea an electric field E, it is accelerated in the direction
+ + vx − α —
of E. On the other hand, if a charged particle
+ + d is released in static or steady magnetic field,
Va nothing happens. The particle remains at rest. But
if the particle is moving across a magnetic field, it
Electron gun Deflecting x
plates
is accelerated at right angles to its direction of
Screen motion.
232 | Physical Sciences

The force on a current element of length dL Gyro Frequency

in a magnetic field is given in vector form by
— —
The number of revolution per second of the
dE = (I × B)dL (N Motor equation) particle in a circular path is the gyro frequency of
I → current, B → magnetic flux density (T), the particle. Introducing v = wR = 2πR gives the
dL → element length (m). frequency
Also (Motor equation applies to moving
charged particles) f=
v
=
eB
2πR 2πm
(Hz) [f = 2πv mveB
q
I = (A) If the particle had an initial velocity com-
t —
Current times length is then equivalent to the ponent parallel to B as well as perpendicular to
—
charge times its velocity, or B, the particle would move in a helical path
—
q with the axis of the helix parallel to B.
IL = L
t If in addition to the magnetic field, there is
L —
= q = qv (Am) also an electric field E, the total force on the
t particle is
Now IdL = dqv — — — — — —
dq → charge in a length dL of a conductor or F = eE + e(v × B) = e[E + (v × B)
a beam. —
— — — E can exert a force on a stationary (or
So dF = d·q(V × B) moving) charged particle and impart energy to it
For a single charge through acceleration. However, there is no force
— — on a charged particle in a static magnetic field
F = e( v × B) Lorentz force unless the particle is moving with a velocity v
— — — — —
⇒ m a = e(v × B) …(1) across the magnetic field B ( v × B ≠ 0). Since the
— e — force is normal to the particle velocity, t h e
a = (v × B) …(2) magnetic field cannot change the particles
m
According to equation (2), the acceleration is kinetic energy. If the particle moves in the same
— — —
normal to the plane containing the particle path direction as B, v × B = 0 and the magnetic field
—
and B. If the direction of the particle path has no effect at all (exerts no force).
—
(indicated by v) is normal to B, the acceleration Cathode-Ray Tube—Magnetic Deflec-
is max. If the particle is at rest, the field has no tion
effect. Likewise if the particle path is in the same
—
direction as B, there is no effect, the particle Let us consider what happens when an elec-
—
continuing unreflecting. tron with velocity | v | = v x enters a magnetic field
—
Particle Radius B (out of page), at a distance x from a screen.
Let us determine the radius R of the circle. vy v
— Screen
The magnitude of the force F (radically inward) vx
− y
on the particle is, vx
— α
F = ma
Electron
= eVB path
This is balanced by the centrifugal force B (out)

mv 2
F =
R
mv 2 It is acted on by a force
or, = eVB — — —
R F = e(v × B)
mv Which accelerates the charge upward in the
or, R =
eB figure (Since e is negative).
 Physical Sciences | 233

As an approx., the acceleration can be taken Potentials and Fields

in the y-direction if the path L is small. Then the
Scalar and Vector Potential—
velocity component in the y-direction is—
L A and φ (or V)
vy = ayt = ay Magnetic and electric fields are given by
vx
— —
F = ma = e( v × B)
— —
B = ∇×A
evxB L eBL
vy =
m vx
=
m — ∂A
and E = –∇φ –
v ∂t
eBL
α = tan–1 = tan–1
y
So —
vx mv x — ∂A
or, E 2 = –∇V –
∂t
vx =

√ 2eV a
m
va = accelerating voltage, V
Q. 1. Find the charge and current distri-
bution that would give rise to the potential.
Thus the deflection angle can be expressed V = 0
 e 
µ0 k
α = tan–1  BL

√ 2mVa A = 
 (Ct – | x |)2^z for | x | < Ct
 4C
0 for | x | > Ct
Compare with α = tan–1 () y
x Where k is a constant, and


√
–
+e 1
C =
±y = xBL
2mVa 
√ ε0µ0
For an electron (e minus) Solution : First we will determine the electric
xBL and magnetic fields,
y = 2·96 × 105 —
√Va
 —
E = –
∂A
√Va (T) ∂t
B = 3·38 × 10–6 µ0 k
xL = (Ct – | x |)^z
2
— —
Rotary Motor or Generator B = ∇×A
If a rotatable rectangular loop or multiturn µ k ∂
— = – 0 (Ct – | x |)2 y^
coil is placed in a uniform magnetic field B and 4C ∂x
equipped with a two-segment commutator µ k
= ± 0 (Ct – | x |)^y
2C
L Ez By
I µ0 kt
I 2
B
− ct ct ct ct
F = (I × B) L x x
− µ0 kt
− µ0 kt
2
2
T = IABN sin θ Nm Loop torque — —
∇· E = 0
—
Average torque ∇· B = 0
Tau =
2IABN
(Nm) — µ k
π ∇ × E = –+ 0 ^y
2
Average d.c. emf generated is — µ0 k ^
2w BAN ∇×B = – z
V = (V) 2C
π
234 | Physical Sciences

–
∂E µ0 kC ^ [In gauge transformation
∂t
= –
2
z A′ → A + ∇λ
∂B µ0 k ^ Equations (i) and (ii) remains invariant)
= ± y
∂t 2 V′ → V –
dλ
Maxwell equations are all satisfied with ρ ∂t
and J both zero. So equation
—
Notice, however that B has a discontinuity at – ∂2 A
∇2 A – µ0 ε0 2 = –µ0 J
x = 0, and this signals the presence of a surface ∂t
– ∂ 2V ρ
∇2 V – µ0 ε0 2 = –
current k in the yz plane, boundary condition ∂t ε0
(σ)
In static field, it reduces to
1 4 1
B – B2 4 = σ × ^n ρ
µ1 µ2 ∇2 V = –
^ ε0
= (k × n)
f –
– ∇2 A = –µ0 J
k + y^ = k × n^ (kc × ^n)
So
–
k = ktz^
Solutions V(r) =
1
4πε0
ρ(r1 )
r
dz ∫
µ
∫
Coulomb gauge –
J(r′)
– A( –r ) = 0 dz′
∇· A = 0 4π r
Important result Consider
–– 1 – – r
φ = E· r and A = (B × r )
2 r
A
Lorentz gauge r′
∂φ ∂V
∇·A = – µ0 ε0 or ∇·A = – µ0 ε0 In the non-static case, this is not the status of
∂t ∂t
the source now that matter, buts rather its
Retarded potentials condition at same earlier time tr (called retarded
– ∂A time) when the “message left”. Since this
E = –∇V –
∂t r
message must travel a distance r, the delay is
– ρ c
(
Substitute into ∇· E =
ε0 ) tr = t –
r
∂ 1 C
∇2 V + (∇·A) = – ρ …(i)
∂t ε0
B = ∇×A
V =
1
4πε0 r
∫
ρ(r′tr)
dτ′
∂E
Mean ∇ × B = µ0 J + µ0ε0
∂t
V(r, t) =
1
ρx′‚ t –
∫
(r
c
dτ′
)
∂ ∂A
∇ × (∇ × A) = µ0 J + µ0ε0
∂t (
–∇V –
∂t ) 4πε0 r

∂V A( –r ‚ t) = 0 ∫
µ J(r–‚ (t – r/c))dτ′
∇·(∇·A) – ∇2A = µ0 J – µ0 ε0 ∇
∂t( ) 4π r
These are called retard potentials
∂2 A
– µ0ε0 2 ρ(r′, ta )
∂ 2A
∂t
∂V
Va (r, t) =
1
4πε0 ∫
r
dτ′
( ) (
So ∇2 A – µ0 ε0 2 – ∇ ∇·A + µ0 ε0
∂t ∂t ) µ J(r–′‚ ta )
= –µ0 J …(ii) Aa (r, t) = 0
4π ∫r
dτ′
 Physical Sciences | 235

In which charge and current densities are

evaluated at advanced time ire
edw
r uc
t
ta = t + nd
C Co

ire
Transmission Lines

w
g
tin
uc
Generally, the wave propagation is divided

into the following two categories.
1. Unbounded media propagation.
2. Bounded media or guided wave propaga-
tion.
(1) In unbounded media propagation, the
propagation is unguided. The uniform plane wave
exists throughout all the space. Therefore, the
energy spreads over a wide area. This type of
propagation is used in radio or T.V. broadcasting.
(2) In bounded media or guided wave propa-
gation, the propagation of plane waves takes place
between boundaries. The boundaries are provided
by the conductors. This types of propagation is
used in telephones, transmission of electrical
energy from generating station to load centres,
connection between a transmitter to an antenna,
connection between computers in networks, etc.
There are two kinds of wave guiding systems.
(a) Transmission lines, and (b) wave guides.
A transmission line consists of two or more
parallel conductors, whereas a wave guide is
generally made up of one conductor.
Transmission lines are used to transmit elec-
tric energy from one point to another point.
Another Definition
Transmission lines is a device, with minimum
energy losses, used to transfer the energy or signal
or power, from a source to a load. (e.g., trans-
mitter to receiver).
We shall consider the transverse electro-
magnetic (TEM) waves guided by transmission
– —
lines. In TEM waves, E and H are both perpendi-
cular to each other and both are transverse to the
direction of propagation. The properties of TEM
waves are same as those of uniform plane waves.
Types of Transmission Lines
Generally the following two types of
transmission lines are used—
(1) Two wires transmission line or parallel
wire line—Two wires transmission line consists
of a pair of parallel conducting wires separated by
a uniform distance as shown in fig. (a)
nd
Co
Fig. (a)
The other type is a shielded pair line. In this
type of transmission line, two parallel wires are
surrounded by an insulator contained in a
cylindrical sheet as shown in fig. (b)
Conducting
wires
Cylindrical
sheet
Insulator
Fig. (b)
These are used in power systems or
telephones. The parallel wire lines are used at
frequencies of the order of 100 mega cycles per
second. Above 200 maga cycles per second, there
is a large loss of energy due to radiation. So, these
lines are unsuitable.
(2) Co-axial transmission line—The coaxial
transmission lines is generally used as high
frequencies. It consists of two concentric conduc-
ters which are separated by an insulating material
as shown in Fig. (c).
Coaxial outer
conductor
Insulating
material
Inner
conductor
Fig. (c)
The outer conductor forms a shield. It
confines the wave between the conductors.
In this transmission line, the electric and
magnetic fields are only confined to dielectric
region. The energy loss due to radiation is quite
small. The losses in dielectric medium increases
with frequencies.
Therefore coaxial lines cannot be used at
frequencies above 1 GHz.
236 | Physical Sciences

Equivalent Circuit Harmonic Time Dependent

The transmission line consists of conductors. Voltage and Current Equations
So, there exists resistance (R) throughout its Considering voltage and current to the har-
length. monic time dependent, the instantaneous line
The resistance R is defined as loop resistance voltage and current are given by
V(x, t) = Re[V0(x)eiωt] V0(x) includes the
per unit length of the line.
When a current passes through a conductor it
}
I(x, t) = Re[I0 (x)ejωt] factor e–(α + jB)x
creates magnetic flux. The flux linkage gives ∂V
Now = jωV
effect of inductance. ∂x
Therefore, there exists a distribution ∂I
And = jωI
inductance (L) throughout the line. [Inductance ∂t
defined as loop inductance per unit length of the Substitute (A) and (B), we have
line]
∂V
In the case of transmission line, the dielectric – = RI + L(jωI)
∂x
material of insulator is not perfect. As a result,
∂V
there is a flow of leakage current i.e., same ⇒ – = (R + jωL)I …(1)
leakage current passes through dielectric material. ∂x
This produces a lekage conductance. Similarly
(a) It is defined as shunt conductance ∂I
– = (G + jωC)V …(2)
between two conductors. dx
In a transmission line, two conductors are Take the differential of equation (5), we get
separated by a dielectric medium. So, there exists ∂2 V dI
a capacitance (C) throughout the line and the – 2 = (R + jω L)
dx dx
capacitance is defined as shunt capacitance
d2 V
between two conductors per unit length of the ⇒ = (R + jω L) (G + jωL)V …(3)
line. dx2
Similarly
So R, L, G and C are known as primary line
constants. As these parameters are distributed d2 I
= (R + jω L) (G + jωC)I …(4)
along the entire length of transmission line, the dx2
line is known as uniform transmission line. Now,
General equivalent circuit of transmission line R + jωL = Z
L R L R L = Series impendance of the line
R
per unit length (Ohm/m)
G + jωC = Y
Output
Input

G C G C = Shunt admittance of the line

per unit (mho/m)
Further γ2 = (R + jω L) (G + jωC)
or, γ = √

(R + jω L) (G + jωC)
Transmission Line Equations = (α + jβ)
Voltage equation— γ → Propagation Constt
∂V ∂I Gives information about the variation of
– = RI + L …(A)
∂x ∂t voltage and current w.r.t. distance along the line.
Current equation— α = Alternation Constt
∂I ∂V This determines the reduction in voltage or
– = GV + C …(B)
∂x ∂t current along the line.
 Physical Sciences | 237

When α is multiplied by the length of line, So the characteristic impedance depends on

this gives total attenuation. R, L, G, C and frequency of operation.
β = Phase Constt Z0 → can also be expressed as
Gives variation of phase position of voltage Z0 = R 0 + jX0
and current along the line.
[R0 and X 0 are real and imaginary part of Z0 ]
The transmission line equations for infinite
transmission line expressed as— The transmission line can be classified in the
following two categories—
–
dV
= ZI  d2 V 2
=γ V (1) Lossless line

dx dx2 …(I)
(2) Distortionless line.
dZ
– = YV
dx
 d2 I
dx2
= γ2 I …(II) (1) Characteristic Impedance of Lossless
Line—A transmission line is said to be lossless, if
Solution of Transmission line Equations— at radio frequencies, the inductive reactance (jωL)
The solution of equations I and II can be is much larger than the series resistance R and
written as— capacitive reactance (jωC) is also much larger
than shunt-conductance G. Thus, in case of loss-
V(x) = V+(k) + V– (x) less line
or, V(x) = V0+e–γx + V0 – e +γx R = 0
V0+ → complex amplitude and G = 0
in +ve x-direction
 – Equivalent circuit of lossless line
V0 → complex amptitude L L L
in-X direction
⇒ V(x, t) = V0 +e–αx

cos(ωt – βx) + V0– e–αx cos(ω t + βx)

Output
Input

C C C
Similarly
I(x, t) = I0+e–αx cos(ω t – βx)
+ I0– e–αx cos(ω t + βx)
Characteristic Impedance
Z0 of line is defined as the ratio of positively
So Z0 =
√jωL
 jωC
= R + jX
0 0

 LC and X = 0
√
travelling voltage wave (V 0 +) to current wave (I0 +)
at any point of the line. R0 = 0

V+ V– Thus characteristic impedance is a real

Z0 = 0+ = 0–
I0 I0 quantity. Loss-less line is pure resistive.
dV (2) Distortionless line—A distortionless line
– = (R + jω L) I
dx is one in which attenuation constant α is fre-
Consider positively travelling wave, we quency independent while the phase constant β is
have V = V0+e–γx linearly dependent on frequency. For a distor-
tionless line,
and I = I0+e–γx
R G
Substituting we gets =
L C
R + jωL

√

∴ Z0 = R + jωL
r M0 =
(R + jωL) G + jωC
=

√ j ωL
R (1 +
R )
(R + jωL) (G +jωC)


√
=
R + jωL j ωC
G(1 +
G)
Z0 =
G + jωC
238 | Physical Sciences

( j ωG C) L
√
C
R 1+ ∴ α = R
=
j ωC
G(1 +
G )
β = ω√ LC
α does not depend upon frequency but β is a

G= √
√ C
R L linear function of frequency.
=
Phase velocity—
ω ω
√C
L 1
Further Z0 = R ≠ jX Vp = =
β ω√
=
0 0
LC √LC

 LC
√
Input Impedance of Transmission Line
R0 =
(l − x)
Transmission line Output
and X0 = 0 Zg
IL
Thus, distortionless line has a real charac- +
teristic impedance like lossless line. It is Vg ~ vi = zin zin ZL VL
important to mention here that a lossless line is −
also distortionless line but a distortionless line is
not necessarly lossless. x = 0 input x = l output
In power transmission, a lossless line is
desired while in case of telephone lines, a Z0 → characteristic impedance
distortionless line is required. γ → propagation constant
Propagation or Constant and Phase Now, In general,
Velocity V(x)
Z = …(1)
I(x)
γ → √
(R + jω L) (G + jωC) is defined as
propagation constant. In general, transmission line equations for
voltage and current are
(1) Lossless line—In case of lossless line,
R = 0 and G = 0 V(x) = V0+e–γx + V0 – eγx
I(x) = I0+ e–γx – I0– eγx
γ = jω√
LC
V+ V–
γ = α + jβ I(x) = 0 e–γx – 0 eγx
Z0 Z0
= jω√
LC (1) Voltage and current at input (x = 0)
⇒ α = 0 Substitute (x = 0)
∴
β = ω√
LC Vi = V0+ + V0 – [ V(x) = Vi]
ω 1 Ii =
1
[V + – V0– ]
Vp = =
β √ LC Z0 0
(2) Voltage and Current at output side
(2) Distortion less line—The conduction for (x = L)
distortionless line is
VL = V0+e–γL + V0 – eγL
R G
= 1
L C IL = [V +e–γL – V0– eγL]
Z0 0
γ = α + jβ
Solve for V 0 + and V 0 –
=

√ ( )
(R + jωL)
RL
C
+ jωC 1
V0+ = [VL + Z0 IL]eγL
2
…(A)

√ CL (R + jωL)
1
= V0– = [VL – Z0IL]e–γL …(B)
2
 Physical Sciences | 239

(3) Voltage and Current at any Point on and tanh (rl) = j tan βl
the Line  ZL + jZ0 tan βl
Zin = Z0 
 Z0 + jZL tan βl
IL
V(x) = [(ZL + Z0 )eγ(l – x)
2
β l → electrical length of the
+ (Z L – Z0)e–γ(l – x)] …(C)
line in degree or radian
IL
I(x) = (b) For a short circuited line ZL = 0
2Z0
 jZ0 tan βl
[(ZL + Z0 )eγ(l – x) – (ZL – Z0)e–γ(l – x)] So, we get[Zin]SC = Z0 
VL  Z0 
Where ZL =
IL = jZ0 tan βl
For lossless line; load power = input power. (c) For open circuited line, Z L = ∞
We consider that
Z 1 + j Z tan βl
  Z 
0
ZL = Z0 L

Z 
L
So Vi = ILZ0eγl [Zin]OC =
Z ZZ + j tan βl
0
0
Ii = ILeγl
Voltage and current at any point is given by
  L
L 
V(x) = (I LZ0eγl)e–γx 1 + j Z tan βl
 Z 
0

Z + jβ tan βl
= Vie–γx = Z0
L

I(x) = (I Leγ l)e–γx

Z 
0
and
L
= Iie–γx
∴ When a finite transmission line is termi-
nated with its own characteristic impedance, there
= Z0 [0 +1j+tan0 βl]
will be no reflected wave. = –j Z0 cot βl
(4) Input Impedance In case of open circuited line (ZL = ∞)
At generator end, the input impedance is (Zin)OC = –jZ0 cot βl
given by—
Vi V0 + + V 0 – In short circuited line (ZL = 0)
Zin = = (Zin)SC = j Z0 tan βl
Ii 1
[V + – V0– ]
Z0 0 (Zin)OC (Zin)SC = Z0 2
Z [V + + V0 – ] Z0 = √
= 0 +0 or, 
(Zin)OC (Zin)SC
[V 0 – V0– ]
(5) Power Consideration
Substitute the values of V0 + and V0 – from (A)
The average power delivered by generator to
and (B) in solve, we get the input terminals of the line is
 ZL + Z0 tanh(γL)
Zin = Z0 
1
 Z0 + ZL tanh(γl) 
(PaV)i = Re[ViIi]x = 0
2
eγl + e–γl The average power delivered to the load is
used = cosh γl 1
2 (PaV)L = Re[VLIL*]x = L
2
eγl – e–γl
and
2
= sinh (γl) 1VL2 RL
2ZL 
=
Special cases—
(a) For a lossless line— 1
= | I | 2 RL
γ = jβ 2 L
240 | Physical Sciences

Transmission Line as Circuit Elements— So Zin = jR0 tan βl

Under normal frequencies, the transmission line = jxin
acts as a communication channel between Here βl can vary from –∞ to +∞.
transmitter and reciver.
∴ Input impedance may be purely inductive
At very high frequency [300mHz to 3GHz], it or purely capacitive depending upon the value of
also acts as ciruit elements like inductor and βl.
capacitor. Ptot :
Considering a lossless line and substituting

Xin = − R0 tan βl
Z0 = R0, the input impedance
 ZL + j R0 tan βl
Zin = R 0 
 R0 + jZL tan βl  π
2
π 3π
2
2π

Special cases—
(1) Open circuited line (ZL = ∞)
ZL → ∞ For very short line i.e., tan βl = βl
(Zin)OC = –jR0 cot βl (Zin)SC = jR0 tan (βl)
= jR0βl
This equation shows that input impedance of


√
an open circuited lossless line is purely reactive. L
= jR0 × ω√
LCl
It can be capacitive or inductive depending C
upon the values of cos βl, when it is +ve or –ve = l(jωL)
respectively plot of input resistance X in = – R0 cot Showing inductive reactance.
βl. (3) Half Wave Transmission Line
For very short length of line βl << 1 λ
l=
2 ( )
λ
Xn = − R0 cot βl

When the length of line is l = n

2 ()
π π 3π 2π nλ
2 2 βl = × = nπ
λ 2
2π Now tan βl = 0
Input reactance of open circuited transmission line ∴ (Zi)λ/2 = ZL
i.e., tan βl ≈ βl λ
jR0 (4) Quarter Wave Transmission line l =
Zin = – 4
tan βl λ
R When length of line is an odd multiple of
= –j0 0 4
βl i.e.,
π

√
L l = (2n – 1)
4
C
= –j 2π λ
l(ω√LC) βl = (2n – 1)
λ 4
1 π
Zin = = (2n – 1)
l(jωC) 2
1 π
= –j tan βl → ± ∞ at tan (2n – 1) ]
ωCl 2
This is inpedance of a capacitor. R0 2
∴ (ZL)λ/4 =
(2) Short circuited line ZL
(ZL = 0) Inverse of load impedance multiplied by
ZL ≈ 0 square at Characteristic-Impedance.
 Physical Sciences | 241

Reflection Coefficient and Standing  1 + l 

Wave Ratio : ZL = Z0  
 1 – l 
Consider the case that power is applied by the
generator at the input end of transmission line. Reverse
There are two possibilities—
ZL – Z0
(1) The transmission line is of infinite length TL =
or terminated by its characteristic impedance. ZL + Z0
Under this situation, all power is absorbed by the Reflection coefficient at the load end.
load at far end. There is no reflected wave. Similarly current reflection coefficient can be
(2) The T.L. is of finite length or terminated obtained.
by the impedance other than its characteristic At another point x where
impedance. In this case, some power is absorbed x′′ = l – x, the voltage reflection coefficient is
at the load end and the remaining power is
given by—
reflected back to the generator.
V – – eγx
Therefore, reflected waves are seen up. There Γx′ = 0 + –γx
is a loss of power due to reflection. V0 e
The incident and reflected wave combine to V0– eγle–γx′
=
produce standing waves of current and voltage V0+e–γl e+ γx′
along the T.L.
The resultant voltage (or current) at any point Γx′ = ΓL e–2γx′
along the line will be a phasor sum of indidence Special Cases—
and wire reflected wave components.
At a certain point two waves interfere in such (1) If the line is terminated at characteristic
way resultant voltage or currents becomes impedance, Z1 = Z0 , then
maximum. Z – Z0
ΓL = 0 =0
Z0 + Z0
No reflected wave.
(2) If the line is short circuited at the receiv-
ing end, ZL = 0, then
0 – Z0
λ λ λ
ΓL = = –1
3λ 0 + Z0
4 2 4
This shows that reflection coefficient is nega-
Reflected voltage at load
TL → tive i.e., reflected wave is reversed. This shows
Incident voltage at load that reflected wave is identical to incident wave
Reflected current at load but travelling in opposite direction.
=
Incident current at load (3) If the line is open circuited at the receiv-
V0– eγl ing end, ZL = ∞, then
TL =
V0+ e–γl
 Z0 
● VL = V0+e–γl + V0 – eγl 1– 
1  ZL 1 – 0
IL = [V +e–γl – V0– eγl] ΓL = = =1
Z0 0  Z0  1 + 0
 V +e–γl + V0 – eγl 1+ 
VL
= Z0  0+ –γl
 ZL
 V0 e – V0– eγl
ZL =
IL This shows the reflected wave is identical to
 V +e–γl incident wave and sign.
= Z0  0+ –γl
 V0 e  (4) If the load impedance is greater or less
than characteristic impedance Z0 , then reflection
 1 + (V0 – eγl/V0 +e–γl)  coefficient is always less than one
 1 – (V – eγl/V +e–γl) = Z
0 0 ZL = 2 Z 0
242 | Physical Sciences

2Z0 – Z0 For a matched line, SWR = 1

Then ΓL =
2Z0 + Z0 ⇒ For a mismatched line SWR > 1
Z0 ∴ Higher is the value of SWR, the greater is
= mismatch between load and line.
3Z0
A High SWR on the line is undesirable
= 0·33 because, it causes a larger power loss and noise.
Z0 Vi & V r → rms value of incident and reflec-
If ZL = ted voltage.
3
| Vmax |

Then ΓL =
(Z3 – Z )
0
0
S =
| Vmin |
| V i | + | Vr |
(Z3 + Z )
0
0
=
| V i | – | Vr |
1 + | Vr /Vi |
= – 0·5 S =
1 – | Vr /Vi |
Maximum reflection occurs when either the
line is short circuited or open circuited. 1 + | ΓL |
=
1 – | ΓL |
Standing Wave Ratio (SWR) ⇒ Reverse ration or
Now we know that if a line is terminated by a S–1
| ΓL | =
load other than Ch. impedance, a reflected wave S+1
exists on the line. The superposition of incident
and reflected wave give rise Standing wave or In case of real (resistive load), we can also
non-progressive wave as shown. express SWR in terms of load impedance ZL, and
characteristic Impedance Z0.
The ratio of max. to min. magnitude of
voltage on the line is called as voltage standing 1 + | ΓL |
SWR =
wave ratio. i.e., 1 – | ΓL |
 ZL – Z0
1+ 
 ZL + Z0
=
 ZL – Z0
Voltage
Vmax 1 – 
 ZL + Z0
Vmin
ZL
Distance =
Z0
Z
S =
Max. magnitude of voltage ∴ SWR = L
Min. magnitude of voltage Z0
| Vmax | Impedance Matching
=
| Vmin | Matched and Unmatched T.L.—Matched
Similarly current, SWR when the load impedance is equal to characteris-
tic impedance
Max. magnitude of current
I = Zin = Z0
Min. magnitude of current
[I ] and reflection coefficient
= max ΓL = 0
[I min]
Significance of SWR is that, it gives a mea- SWR(s) = 1
sure of mismatch between load and transmission If this condition is not satisfied, the line is
line. said to be mismatched.
 Physical Sciences | 243

Impedance matching by quarter wave or, Z0 ′ = Zin Z L

√
transformer—The mismatching between load ZL ∴
and Ch. impedance Z0 can be removed by = 
√ Z0 ZL [ Zin = Z0 ]
λ Single stub Matching—The impedance
inserting a line of length with Ch. impedance Z0
4 matching can also be achieved by connecting
λ sections of open or short circuited lines called
between original transmission line of length
4 stub or tuning stub in shunt with the main line at
with Ch. impedance Z 0 ′ between original some point or pouts. This method of achieving is
transmission line and load as shown is figure. called stub matching.
Original Transmission Line Basically of two types.
(i) Single stub matching
Inserted λ/4 line
(ii) Double stub matching.
Z0 Zin ZL Consider single stub matching.
In this method, a short circuited stub of
Quarter wave length l is connected in parallel with the main line
transformer at some distance of from the load.
λ/4 l
− jB Yin = Y0 + jB
YS =
λ
The transmission line is called as quarter
4 Y0 ZL
wave transformer because line acts as a trans-
former.
Calculate t h e C h . impedance Z 0 ′ of the
line. d
For a lossless line, the input impedance as ob-
Short circuited stub
served by looking into quarter wave transformer is
 ZL + j Z0 ′ tan β l
Zin = Z0 ′ 
The short circuited stub are preferred as
 Z0 ′ + j ZL tan β l
compared to open circuited stub because open
circuited stub radiates same energy at high
λ λ frequency.
For a line, l =
4 4 1
2π λ Let, the Ch. admittance Y0 = Admittance
βl = × Z0
λ 4 of stub YS = – YB
π B → Susceptance of the stub.
=
2
Line input at AB = Yin = Y0 + j B → imaginary
π
tan βl = tan = ∞ ↓ part
2
Now, Zin = Z0 ′ Real part

( Z ) + Z ′ 
The principle of this technique is to connect a

 j tan βl 
L
0
short circuited stub of appropriate length l and at a

( Z ′ ) + Z 
= certain distance d from load end in parallel such
that its input admittance is equal to the
 j tan βl 
0
L negative of imaginary part of line input
admittance to the right of stub.
 0 + Z0 ′
Z0 ′ =  0 + ZL  Now, the line admittance towards the left of
the stub
(Z0′)2 = (Y 0 + jB) – jB
=
ZL = Y0 [∴ match is achieved]
244 | Physical Sciences

Theoretically, we can add inductors or capa- ZL – Z0

citors in parallel with the T.L. but these elements ΓL =
ZL + Z0
are usually too lossy at the desired frequency. So, = (ΓL)r + j(ΓL)i
their uses is avoided.
↓ ↓
Double Stub Matching Real Imaginary
part part
A single stub cannot be placed physically in
an ideal location. So, “double stub technique is The normalized impedance Zl is given by
employed. In this method, two short circuited ZL
parallel stubs are attached to the T.L. at fixed Zl = = r + jx
Z0
positions and the lengths are adjusted in such a
way to provide perfect matching. The spacing From equation (1)
between two stubs is usually (λ/8, 3λ/8, 5λ/8 etc).  ZL
Z  – 1
3λ
The separation of is an optimum.  0
8 ΓL =
 ZL
Z  + 1
ls2 ls1
 0
Zl – 1
=
Zl + 1
ZL
(r + jx) – 1
or, (Γ) r + j(Γ)i =
O (r + jx) + 1
B Solving, we get
l2 l1 (1 + Γr) + jΓi
r + jx =
(1 – Γr) – j Γi
Smith Chart—Graphical method to solve the Rationalize
problems of T.L. 1 + Γr 2 – Γ i 2 2Γi
r + jx = +j
It shows, in a very simple way, the SWR (1 – Γr)2 + Γi2 (1 – Γr)2 + Γi2
corresponding to as given lead impedance also Comparing the real and imaginary parts on
line impedance at any desired points as one moves both sides, we get
along the line.
1 – Γr 2 – Γ i 2
It consists of a group of concentric circles r =
(const resistance and conductance circles). (1 – Γr)2 + Γi2
On this chart the normalized values of 2Γi
and x =
impedance ZL to characteristic impedance Z0 (i.e., (1 – Γr)2 + Γi2
Z 1 Rearranging the equation of r, we get
Z1 = L) and normalized admittance YL = .
Z0 ZL
Construction—The Smith chart consists [ Γr –
r
(r + 1) ]2
+ Γi 2 = ( )
1
1+r
2

of two sets of perpendicular circles. One

is identical with
Z  R 
set indicates the ratios  L or  L; where (x – g)2 + (y – h)2 = a2
 0
Z  Z0  rad. a, centre at (g, h).
R L → resistive component of line impedance and
Compare, we get centre of constt resistance
Z0 → characteristic impedance which is also
resistive for loss-less line. The second set of circle
gives the ratio jX/Z0 , where X is reactive compo-
circle is ( 1+r
r
)
‚ 0 ‚ with radius
1
(1 + r)
.

nent of the line impedance. Similarly, arranging the equation of x, we get

First of all, determine the equation of these
circles. The reflection coefficient of a T.L. (
(Γr – 1)2 + Γi –
1 2
x ) ()
=
1
x2

Constant reactance circle with centre (1‚ x1)
1 (v) Optical fibre, etc.
rad .
x Essential feature of wave guide propagation
R=∞
2 1.5 1
R=0
Constant (r) circle
jx
∞ k=0
− jx
SWR circle—This is the circle with centre as
the centre of the chart and the rad. of the circle is
the distance between the centre of chart and the
normalized impedance.
SWR of the design—The point of intersec-
tion of SWR circle and two axis of the chart (rhs
of the centre of the chart) gives the SWR of the
design.
Load line—The line passing through the
centre of the chart and the given normalized
impedance is known as load line.
Application of Smith Charts
(1) To find the location of max voltage and
minimum voltage from the load.
(2) To find the SWR.
(3) To find unknown load impedance from
reflection coefficient.
(4) To find position and length of the stub
used in impedance matching.
Wave Guides
A hollow conducting, metallic tube of
uniform cross-section usually filled with air, for
transmitting electromagnetic wave by successive
reflections from inner walls of the tube is called a
wave guide.
Wave guides are of different forms such as—
(i) Parallel-plate guide
(ii) Rectangular wave guide
Physical Sciences | 245
(iii) Cylindrical wave guide
(iv) Dielectric wave guide
is that it exhibits cut off characteristic frequency
similar to that of a high pass filter.
The guided electromagnetic wave in a wave
guide can be classified in the following three
categories—
(1) Transverse electric mode (TE)—In this
mode, there is a component of magnetic field
in the direction of propagation and no electric
field component in the direction of propagation,
i.e., the wave will have its electric field compo-
nent only in transverse plane of the direction
of propagation.
(2) Transverse magnetic mode (TM)—In
this mode, there is a component of electric field in
the direction of propagation and no component of
magnetic field in the direction of propagation.
(3) Transverse electromagnetic mode
(TEM)—In this mode, both the fields components
are not in the direction of propagation.
Waves between Parallel Plane
The equation for plane wave propagating in
z-direction varying sinsioidally with the space
enclosed by two conducting planes which are
infinite in extent.
γ ∂Hz
Hx = – 2
h ∂x
jωε ∂Ez
And Hy = – 2
h ∂x
γ ∂E
Ex = – 2 z
h ∂x
jωµ ∂Hz
And Ey =
h2 ∂x
Consider three cases—
(1) Ez = 0 and Hz ≠ 0, this is the case of
transverse electric mode (TE-mode).
(2) H z = 0 and Ez ≠ 0, this is the case of
transverse magnetic code (TM mode).
(3) Ez = 0 and Hz = 0 i.e., both E and H do not
have the components in the direction of propaga-
tion. This is called as transverse electromagnetic
mode (TEM mode).
Cut off frequency (TE mode)
m
fC =
2a √µε
246 | Physical Sciences

Phase velocity (TE mode)

(vg )phase =
v0 Hx0 =
jωεC
h2 (nπb ) sin(mπxa) cos(nπyb)
  λ0  2
h ( a )
cos(
a ) ( b )
jωεC mπ mπx nπy
 1 –  λC  Hy0 = – 2 sin

Group velocity (TE mode) Cut off frequency—

(v g )group
  λ0  2
= (v g )phase  1 –   
fC =
2π
1
√µE
mπ
a [( ) + (nπb ) ]
2 2 1/2

  λC 
Cut off wavelength—
Wavelength within guide (TE mode) 2
λC =

√

1 1 1
λ02
= +
λC2 λg2 [( ) ( ) ]
m 2
a
+
n 2
b
Electric and magnetic lines of TE10 βmin)
Minimum value of phase factor (β
(dominant mode)

√
[ ( )] ωC 2
βmin = ω√
µE
X X 1–
ω
Phase velocity
E v0
Ey
vp =

√

E - Field ωc
z z [ ( )] 1–
ω
2

X Guide wavelength
1 1 1
= +
λ02 λC2 λg2
Wave impedance
η
zg =

√
 [ ( )] ωc 2
z 1–
ω
Magnetic and electric lines of TM 10 mode
For TE mode
X X Field equations (Ez0 = 0)

H Hz0 = C cos (mπxa) cos(nπyb)

C( ) cos(
a ) ( b )
jωµ nπ mπx nπy
Ex0 = + 2 sin
z z h b

C( ) sin(
a )
cos(
b )
jωµ mπ mπx nπy
Rectangular Wave Guide Ey0 = – 2
h a

h (a )
cos(
b ) ( a )
Field equation for TM mode [H z = 0 for TM jβ mπ nπy mπx
Hx0 = 2 C sin
mode]

h (b )
cos(
a ) ( b )
jβ nπ mπx nπy
Ez = C sin ( ) ( )
mπx
a
sin
nπy –j βz
b
e Hy0 = 2 C sin

Cut off wavelength

h ( a)
cos(
a ) ( b )
jβC mπ mπx nπy
Ex0 = – sin
2 v 2
λC = 0 =


√
fc
h (b ) ( a )
cos(
b )
jβC nπ mπx nπy
Ey0 = – 2 sin [(ma) + (nb) ]
2 2
 Physical Sciences | 247

Guide Wavelength The field components for TE waves in circu-

λ0 lar guides
µω
λg =
  λ0  2 Er = H
β φ
 1 –  λC 
µω
Eφ = – H
TM waves in Cylindrical Guides β r
The Field components are
β Ez = 0
Er = H
ωε φ At r = a, the normal field value Hr must be
β zero. This requires that
Eφ = – Hr
ωε
Hz = 0 Jn ′(ha) = 0

Solved Questions

∫
r
Q. 1. Find a charge distribution that would q
produce Yukawa potential ⇒ Q = – e –r/a r dr
a2 0
q e–r/a
φ =
4ππε 0 r ⇒ Q = q [(ra + 1) e –r/a
]
–1 …(3)
also state why the total charge be zero ? If r → 0, Q = –q.
Solution : Poisson’s equation is But from the equation of φ
ρ
∇2 φ = – . r → 0, φ → ∞.
ε0
The potential φ, is spherically symmetric so Thus, for the charge at the nucleus the
that singularity exist at r = 0, if Q′ be the point charge
at origin, from Gauss’ law
1 ∂ ∂φ ∂2 φ 2 ∂φ
∇2 φ = 2
r ∂r ( )
r2
∂r
= 2+
∂r r ∂r
∫ E ds =
1
(Q + Q′)
∂2 φ 2 ∂φ ρ ε0
Therefore, + = – …(1)
∂r2 r ∂r ε0
q e–r/a
[
or ε0 E r (4πr2 ) = Q′ + q
r
a ( )
+ 1 e–r/a – 1 …(4) ]
Again φ = …(2) [from (3)]
4πε0 r
∂φ
Differentiating equation (2) twice, we get
∂2 φ
Also Er = – =
∂r 4πε0
a –r/a 1 1
e –
ar r2[ ]
∂r2
=
q –r/a 2
4πε0
e [ +
2
+
r3 ar2 a2 r
1
] So from (4)
∂φ q e–r/a q e–r/a Q′ = +q
= – –
∂r 4πε0 ar 4πε0 r2 Therefore, Yukawa potential arises due to a
Putting these values in equation (1) point charge +Q at origin and an equal and
q e–r/a opposite charge –Q surrounded. The net charge
ρ = – (Q + Q′) is zero.
4πa2 r
which is the desired charge density. Q. 2. Show that the potential at a distance
Amount of charge enclosed within a spherical r1 and r 2 respectively from centres of a long
shell of radii r and (r + dr) is given by parallel pair of wire of negligible cross-section
dQ = (4πr2 dr)ρ and having equal and opposite linear charge
density λ coulomb metre is given by
∫
r
q e–r/a
Q = – × (4πr2 dr)
0 4πa2 r λ r
φ = log 2.
(Putting the value of ρ) πε
4π 0 r1
 Unit-4

QUANTUM PHYSICS AND APPLICATIONS

Postulates of Quantum Mechanics 2
Operator for Px
(1) First postulate—Information about the ∂
properties of a system of particles is contained in a
function of all co-ordinates and time. This function
2
P x = ih
∂x (– ih ∂x∂ )
∂2
is known as the wave function ψ. The wave = – h2
∂x2
function should be acceptable, i.e., single valued, Operator for total energy (Hamiltonian
continuous and finite everywhere between zero operator)
and infinity. 2 2 2
px + py + pz
(2) Second postulate—Information regarding Since K.E. =
2m
any physical variable of a system can be extracted ^ ∂
and px = – ih etc.
from the wave function with the help of corres- ∂x
^ = – h ∇2 + V
ponding operator for the variable. 2

If f is a physical variable and F is its operator H (x, y, z)

2m
then Fψ = f ψ. The left hand side shows the ∂ 2 ∂ 2 ∂2
operation of F over ψ and right hand side shows where ∇2 ≡ ( +
∂x2 ∂y2 ∂z2
+ )
an ordinary multiplication. Operator for angular momentum
^L = – ih y ∂ – z ∂
e.g., energy operator ih
∂
∂t
x ( ∂z ∂y )
^
∂ ⇒ Lx = – ih (yp z – zpy)
and momentum operator – ih
∂x
Heisenberg’s Uncertainty Principle
(3) Third postulate—The expectation value
of a variable f is This is principle which separates quantum
mechanics from classical physics. According to
∫ ψ* Fψ dτ
〈f〉 = this principle it is impossible to specify the posi-
∫ ψ*ψ dτ tion and momentum of a particle simultaneously
where ψ* is the complex conjugate of ψ and dτ with unlimited accuracy. i.e., there is always an
means integration over all the values of all uncertainty in either or both the specifications.
the space co-ordinates. If ψ is a solution of a And the product of uncertainty is greater than h.
Schrodinger wave equation then a ψ will also be ∆x . ∆px ≥ h
a solution of wave equation, where a is a simple The uncertainty relation between the time
number. So ψ is choosen in such a manner that energy and angular position-angular momentum
denominator of above equation is unity and thus can be given by—
the wave function is said to be normalized ∆E·∆t ≥ h
〈 f 〉 = ∫ ψ* Fψ dτ and ∆θ · ∆J ≥ h
 Physical Sciences | 299

Uncertainty relation is a fundamental law of The divergence of probability current is

our nature and independent of our measurement
techniques.
Schrodinger Equation
This is the basic wave equation which
describes the behaviour of matter waves. The
equation can be written in two forms.
always equal to the time gradient of probability
density.
Degeneracy—If there are two or more
linearly independent wavefunctions belonging to
the same energy state E, the state is said to be
degenerate. For such type ‘g’ functions the state is
known as ‘g’ fold degenerate state.

(i) Time dependent equation— Orthogonal and Normalised Wave

Functions
h2 ∂ψ
– ∇2 ψ + Vψ = i h If the product of a function and the complex
2m ∂t conjugate of another function vanishes on
(In three dimensions) integrating over the interval a ≤ x ≤ b then the
functions are said to be orthogonal in the interval
where V is the potential energy associated with ( a . b). If the functions are ψ1 (x ) and ψ 2 (x )
the de Broglie plane wave. The solution of above respectively then orthogonality condition will
equation will be of the form. be—
–i b
ψ = Ae h
(Et – px)
(One dimensional solution) ∫ a ψ1 (x)·ψ2 * (x) dx = 0
(ii) Time independent equation (Stationary Also we know that the probability of finding
state wave equation)— a particle in the volume element dV, is given by
ψψ* dV or |ψ*|2 dV and the total probability of
h2
⇒ ∇2 ψ + (E – V) ψ = 0 finding a particle in entire space is always unity
2m then
(In three dimensions) ∫ |ψ2 | dV = 1 or ∫ ψψ* dV = 1
The equation is applicable when the total Every wave function satisfying the above
energy of the system is independent of time. equations are termed as normalised wave function.
Normally the ψ is not a normalised wave
A solution of this equation will be function. If ψ is a solution of Schrödinger’s
ψ = Ae – i (ωt – kx) equation the Aψ will also a solution. We choose
(One dimensional solution) the value of A such that it makes the function
normalised.
Probability Current Density ∫ Aψ Aψ* dV = 1
The equation A2 ∫ ψψ* dV = 1
 ih  1
∂ A2 =
(ψψ*) + ∇ – (ψ*∇ψ – ψ∇ψ*) ∫ ψψ* dV
∂t  2m 

is known as equation of continuity. This equation Complementarity Principle

can be obtained by two simple operations of Neils Bohr who was guided by the empirical
complex conjugates on time dependent equation. fact of dual nature of matter, introduced the
The above equation c a n b e written i n t h e complementarity principle which states that the
form of probability current J and probability wave and particle aspects of matter are
density P. complementary rather than being contradictory,
→ dP both equally essential for a full description of the
∇· J + = 0 phenomenon.
dt
A single experiment can display either of
→ ih these nature but both nature of matter cannot be
Where J = – (ψ*∇ψ – ψ∇ψ*) displayed by same experiment simultaneously.
2m
300 | Physical Sciences

Quantum Mechanical Treatment This zero point energy is a consequences of

of Linear Harmonic Oscillator uncertainty principle and does not vanish even at
absolute zero temperature too.
Consider a particle of mass m which oscillate
in linear harmonic motion along X-axis about an Tunneling through a Potential Barrier
equilibrium position. As the particle moves away Tunneling is a purely quantum mechanical
from its mean position a restoring force ^r = – kx phenomenon. Let a particle of mass m is travelling
appears which tends the particle back towards its in + X direction with potential energy E and at a
means position. Here k is known as elastic force point it face a region of high potential energy V.
constant. The oscillation frequency will be This change of P.E. takes place abruptly. First
region lies between – ∞ to 0 and second lies from
√ m
1 k
ν = 0 to + ∞. Needless to say classically the particle
2π cannot be found in region II.
or ω =
 mk
√
The potential energy can be given as
1 2
V(x) = kx
2
This value of P.E. can be substituted in
Schrodinger’s time independent wave equation.
d2 ψ 2m
+ (E – V) ψ = 0
dx2 h 2
∂2 ψ 2m
i.e., +
∂x2 h 2 (E – 12 kx ) ψ = 0
2 Fig. : Classically the probability of finding the particle in
region II is zero until it gains an additional P.E.
equal to (V-E)
This equation has a solution only at the values But quantum mechanics shows an interesting
En = n +
1
2( ) hω (n = 0, 1, 2, …) ability according to it there is always a small
but definite probability of finding the particle
This shows that a quantum mechanical linear in region II, without having any extra potential
oscillator exhibits only a discrete series of energy. This leaking out is known as barrier
energies. Here n is an integer quantum number. penetration or tunneling. This has been confirmed
Unlike classical oscillator, the quantum oscillator through numerous atomic and nuclear particles.
possess zero point energy corresponding to lowest Let us examine its quantum mechanical
quantum number (n = 0) treatment.
1 In region I, the motion of particle is given by
E0 = hω Schrodinger’s time dependent equation.
2
d2 ψ1 2mE
or E0 = hν
1 + h 2 ψ1 = 0
2 dx2
d2 ψ 1
V or + k 21 ψ1 = 0 …(1)
9
– hω E4 dx2
2


√
7 2mE
– hω E3
2 Where k1 =
5
h2
– hω E2
2
ψ1 is the wave function in region I.
3
– hω E1
2 The general solution of equation (1) be
1
– hω
2
E0 ψ1 = A1 exp (ik1x, x) + B 1 exp (– ik – 1k1 x)
x …(2)
0 at x < 0.
 Physical Sciences | 301

Also for region II  2d 

D = D0 exp – h 
√
2m (E – V) 
d2 ψ2 2m  
+ 2 (E – V) ψ2 = 0
dx2 h
where D0 ~
– 1
d2 ψ 2
or + k 22 ψ2 = 0 …(3)
dx2


√
2m
where k2 = h 2 (E – V)

The general solution of equation (3) be

ψ2 = A2 exp (ik2x) + B2 exp (– ik2x) …(4)
at x > 0
A1 exp (i k 1 x , x) corresponding to wave
travelling along +X direction and B1 exp (– ik1x)
corresponding to wave moving in region I in –X
|B |2
direction, i.e., reflected wave. The ratio 1 2 gives
|A1|
the coefficient of reflection ‘R’ from the barrier.
In the same fashion A2 exp (ik2x) corresponds to
wave propagating along +X axis in region II and
B 2 exp (– ik2x) corresponding to reflected) wave
in region II, which do not have any physical
|A |2
existence, B 2 = 0. Again the ratio 2 2 is the
|A1|
penetrability through interface and can be given
by D.
Fig. : Quantum mechanically the particle have a small but
definite probability of existence in region III (It can
leak through higher P.E. region II. In region III it
will possess less energy, i.e., higher wavelength.
Angular Momentum
The moment of linear momentum is known
as angular momentum. For a given particle
angular momentum is given by the cross product
of its position vector with linear momentum.
→ → →
L = r ×p
It is easily shown that

For the barrier height V > E, the wave vector

→
Lx = ypz – zpy 
is imaginary. →
 …(1)
Ly = zpx – xpz

k2 = ik
and
→
Lz = xpy – ypx 

√
2m (V – E) →
Where k = where x, y, z are the components of r and px, py
h
and pz are the components of linear momentum in
will be a real number. The wave function ψ1 and three mutually perpendicular axis.
ψ2 be The quantum mechanical equivalent of set (1)
ψ1 = A1 exp (ik1x) + B1 exp (– ik1x) equations are—
}
ψ2 = A2 exp (– kx)
Since A 2 ≠ 0 there is a probability of
^L = – i h
x [y ∂z∂ – z ∂y∂ ]
penetration to a certain depth y [ ∂x ∂z]
^L = – i h z ∂ – x ∂
…(2)
^L = – i h x ∂ – y ∂ 
 2x
ω ~– |ψ2 | = A22 exp – h 
√

2m (E – V)  z [ ∂y ∂x] 
 
The second set of equations can be written in
For a barrier of finite depth d this probability polar co-ordinate form.
will be
302 | Physical Sciences

→
[sin φ ∂θ∂ + cot θ cos φ ∂φ∂ ] 
In compact form, we can write
Lx = i h J×J = ihJ

L = – i h [– cos φ + cot θ sin φ ]

→ ∂ ∂


…(3) Ladder Operators
y
∂θ ∂φ
→ ∂ Two new operators J+ and J– are known as
L =–ih
z
∂φ ladder operators and can be defined as
Using the second set of definitions it is easy J+ = Jx + iJy
to show that Lx, Ly and Lz satisfies the following and J– = Jx – iJy
commutation rules.
ladder operators do not commute but J 2 commute
[Lx, Ly] = i h Lz with both ladder operators
[Ly, Lz] = i h Lx i.e., [J2, J±] = 0
[Lz, Lx] = i h Ly Also [J2, J±] = ± h J±
In compact form these can be quoted as and [J+ J–] = 2 h Jz
L ×L = ihL ^
Eigen Values of L z
Spin Angular Momentum
To find out the eigen values of ^Lz we need to
The electron also possesses spin motion and satisfy the relation
hence contributes to the total angular momentum.
It is denoted by S. It follows the same commu- Lz ψ = kψ
tation relations as those of orbital angular ∂ψ
or – ih = kψ
momentum. ∂φ
[Sx, S y] = i h Sz The general solution of this equation be
[Sy, S z] = i h Sx ψ = f (r, θ) eikφ/h
(f is a function of γ and θ)
[Sz, S x] = i h Sy An increment of 2π in φ does not change the
and [S2 , S x] = 0 wave function
[S2 , S y] = 0 f (r, θ) eikφ/h = f (r, θ) eik (φ + 2π)/h
[S2 , S z] = 0 or eik·2π/h = 1
2πk
where S 2 = ^S 2x + ^S 2y + ^S 2z or h = 2mπ
The Total Angular Momentum or k = mh
Operators where m is an integer. Thus eigen values of Lz are
The total angular momentum which may given by
include the spin contribution is conveniently Lz = m h
denoted by and the eigen function be
J = (Jx, Jy, Jz) ψeimφ f (r, θ)
The total angular momentum is defined as
generalized angular momentum operator J as any Total Angular Momentum and Pauli
Hermitian operator whose components, satisfy the Spin Matrices
commutation rule The electron in an atom possesses orbital
[Jx, Jy] = i h Jz momentum as well as spin moments. So the total
angular momentum is a vector sum of orbital and
[Jy, Jz] = i h Jx spin moments.
[Jz, Jx] = i h Jy J = L+S
 Physical Sciences | 303

where J and S represents orbital and spin h

momenta. The components of J can be written as Jz =
2 ( 10 –10 )
^ ^ ^
J x = Lx + S x  and J2 = h
3
(10 01)
^ ^ ^
 4

J y = Ly + S y
^ ^ ^ We can also write
and J z = Lz + S z h
J = σ
Their mutual commutation relations will be 2
^ ^ ^
[J x J y] = i h J z
^ ^ ^
Where σx = 1 ( 10)
0

[J y J z] = i h J x
= (
i 0 )
0 –i
σy
^ ^ ^
[J z J x] = i h J y
= (0 –1)
1 0
^ ^ and σz
Also [J x J x] = 0
^ ^
[J y J y] = 0 Here σ x σ y and σ z are known as Pauli spin
matrices. The following commutation takes place
^ ^ among Pauli spin matrices.
[J z J z] = 0
^ ^ σx2 = σy2 = σz2 = 1
and [J x J 2 ] = 0
^ 2 ^2 ^ 2 ^ 2 ^ ^ ^ σx σy = iσz 
J = J x + J y + J z commutes with J x, J y and J z σy σz = iσx 
^ σz σx = iσy 
The eigen values of operator J 2 (choosing a
^ σx σ y + σ y σ x = 0 

unit for which h = 1) are the numbers J (J + 1) σ y σ z + σz σy = 0 

where J = 0, 1/2, 1, 3/2, 2, …… σ z σx + σ x σ z = 0 
^
and each of these is (2J + 1) fold degenerate. Degeneracy of Energy Levels of
^ Hydrogen Atom
The eigen values of J x for a given value of j is
m = J, – J + 1, – J + 2, …… J – 2, J – 1, J. The energy states corresponding to which
Thus it appears that total angular momentum there are more than on quantum states are known
component Jx can have half integer values, on the as degenerate energy levels. And the number of
other hand component of orbital angular moments quantum states represents degree of degeneracy.
can have only integer eigen values. The nature The energy eigen values of hydrogen atoms
provides systems which have half integral values can be given by
of angular moments, i.e., spin of electron. We can
1 2π2 µe4
not write spin angular momentum as an operator En = – ,
(4πε0 ) n2 h2
2
operating on the co-ordinates of electron. Spin has
no classical analogue and is a purely quantum where µ is the reduced mass of H-atom.
number. It is introduced in an adhoc manner in Here energy eigen values depends on principal
non-relativistic theory. Dirac theory is a correct quantum number n only. So the energy levels are
explanation of it. degenerate with respect to l and m.
^ The values of l for a given n (= 0, 1, 2, 3, …)
The total angular momentum matrices for J =
1/2 are will be
h 0 1 l = 0, 1, 2, … (n – 1)
Jx = ( )
2 1 0 and the values of m for a given l be
h 0 –i m = 0, ±1, ±2, ±3, … + ±l
Jy = ( )
2 i 0 So m can have (2l + 1) values.
304 | Physical Sciences

In general form the total number of wave From definite integral identity
function or degree of degeneracy ∞
∫
n!
xn e–αx dx = n + 1
n–1 α
∑
0
= (2l + 1)
l=1 2
Now 4π2 A2 × = 1
n–1 n–1 (2β/r0)3
= 2 ∑ l+ ∑ 1
β2

√
l=0 l=1
or A =
=
2n(n – 1)
+ n = n2 – n + n πr0 3
2 So the normalized wave function is ground
= n2 state be
So the total values of m and l for the given n
are n2 .
For first orbit n = 1
Ψ =

√ β3 –βr/r0
πr0 3
e

or n2 = 1 Time Dependent Perturbation

It is a non-degenerate state whereas all other
states n > 1 are degenerate, e.g., n = 2 is four fold The time dependent perturbation theory
degenerate. (sometimes known as theory of quantum
transitions) is invoked in case of an external
The Wave Function of Deuterium perturbation V(r, t) causing an unpurturbed system
Ground State ^
of Hamiltonian H0 to make a transition from a
Deuterium is a two nucleon system. In which stationary state Ψk0 to another stationary state Ψl0 .
a bound state of single proton and a neutron takes
place. The force of interaction inside the nucleus As we know time-dependent Schrodinger
can be explained in terms of Yukawa potential equation is given by
∂Ψ
V0 ih
V(r) = –
r/r0
exp (– rr )0
∂t
= HΨ
and the Hamiltonian for a system subjected under
where V0 is the peak P.E. and r0 is the range of the time dependent perturbation is
V(r). ^ ^
H = H0 + V(t)
The Schrodinger equation for deuterium be
^
[r1 ∂r∂ (r ∂r∂ ) + V ]Ψ + 2m–h Eψ = 0
2
2
(r)
2
Let the complete set of eigen functions of H0
be Ψ n which are the solution of
0

^
–h 1 ∂ ∂ H0 Ψ 0n = E0n Ψ 0n
2m [r ∂r ( ∂r)
or r 2 + V ]Ψ = Eψ
2
(r)
The wave function representing the actual
The general solution of this equation will be state Ψ(t) can be expended in terms of Ψ0n as
r/r0 0
Ψ = Ae –β Ψ(t) = ∑ bn (t) Ψ0n = ∑ Cn(t) Ψ 0n e–i En t /h
where A is normalization constant and P is a n n
variable. Now, we have
According to normalization condition ∂Ψ ^ ^
∞ (
i h (t) = H0 + V Ψ(t)
∂t
)
∫ 0
Ψ* Ψ dλ = 1
and on introducing Ψ(t) gives
or ∫ (Ae–β r/r0) (Ae–β r/r0) (4πr2 dr) = 1 dC
i h l = Σ Cn (t) Vin(t) exp(i ωlnt)
∞ dt
= 4π 2 A2 ∫ 0
r2 (e–2β r/0)r = 1 where νen = (E0l – E0n )/h
∞
Vin = ∫ Ψ 0*
∫
2!
l V(t) Ψ n dt
0
But r2 e–2β r/r0 dr = and
0 (2β/r0)3
= Matrix element
 Physical Sciences | 305

At t = 0, we take the system to be in a A N particle system can be described by

particular state Ψ 0k so that Ψ (0) = Ψ 0k and thus wave function
C n (0) = δnk . Let us assume that at time t Ψ(1, 2, 3, … r… S… N) where 1, 2, …r …
amplitudes have not changed much from their etc. represent the space and spin co-ordinates of
values at t = 0, i.e., we replace Cu (t) by C u (0) on particle 1, 2, …r… etc.
right hand side of above differential equation for If an operator P is operated on this function,
C e and we get which may or may not leave this function
unchanged i.e.
∫
t
1
C l(t) = V exp (i ωlkt)dt
i h 0 lk PΨ (1, 2, 3 …r…S…N) = ±Ψs(1, 2, 3
exp (i ωlkt) – 1 …r…S…N)
= –Vlk The positive sign corresponds to symmetric
h ωlk
wave function particles known as Boson which
Since at t = 0 the system was in a state Ψ0k , the obeys Bose-Einstein statistics.
probability of transition to a state Ψ 0l at time t is Whereas negative sign corresponds to anti
ωlkt symmetric wave function particles known as
(V
P k → e = | Cl(t) |2 = h
lk) ( )
2 sin
2
2 fermions which obeys Fermi-Dirac statistics.
ωlk 2
( ) 2
Scattering Theory
The differential cross section—
This equation gives the probability of transi-
dσ (θ, φ) 1 dN (θ, φ)
tion to one particular state Ψ 0k . =
dΩ Jinc dΩ
Now the probability of transition to any one Scattering Amplitude of Spineless
state of the group of closely spaced states is given particles—
by
– → –
W = ∫ Pk – l Pl (E0l ) dE0l ψ (r ) = φinc ( r ) + φsc (r )
where Pe(E0l ) dE0l represents the number of final   – –
––
ik · r
e
– A  eik 0 ·(r ) + f (θ‚ φ)
~ 
states in the energy interval dE0l , since Pk – l is r 
sharply peaked around E0l = E0k , we can obtain So DCS is
2π
W = h | Vlk |2, P l (E0k )t dσ 1 dN k
= = |f (θ, ψ)}2 …(1)
dΩ Jinc dΩ k0
The transition per unit time will be
dw 2π Since the normalization factor A does not
= h | Vlk |2 Pl (E 0k ) contribute to the differential cross-section, we will
dt
be taking it equal to one. For elastic scattering k0
This equation is known as Fermi’s golden
is equal to k, hence equation (1) reduces to
rule.
dσ
Identical Particles = |f (θ, φ)|2
dΩ
The identical particles are those which can be
characterized by their invariable mass, charge and Scattering Amplitude—
angular momentum etc. In classical physical µ ^ ψ>
f (θ, φ) = – < φ |v|
identical particles are distinguishable and can be 2πh2
identifical by their fixed trajectories. But in – – ––
–
quantum mechanics the concept of fixed trajecto- where φ (r ) is a plane wave, φ (r ) = e ik ·r , and k is
ries have no meaning. Moreover the definition of the wave vector of the scattered wave. The
distinguishability changes. Here if the position differential cross section is then given by
and spin co-ordinates of two of these particles
dσ µ2
interchanged, we cannot measure the change in = |f (θ, φ)|2 = 2 4 | < φ | ^v | ψ > |2
dΩ 4π h
the system physically.
306 | Physical Sciences

The Born approximation— where σl are called the partial cross sections
The first Born approximation— corresponding to the scattering of particles in
various angular momentum states.
The differential cross section is
The optical theorem is a connection between
µ2 →
∫
dσ ––
= |f (θ, φ)|2 = 2 4 | eiq·r' V ( r' ) d3 r'|2 f (0) and σ :
dr 4π h
4π 4π ∞
– – Im f (0) = σ = ∑ (2l + 1) sin2 δt
where q = k 0 – k k k l=0
In elastic scattering The physical origin of this theorem is the
∞ conservation of the particles (or probability).
∫
2µ
f (θ) = – 2 r' V (r') sin (qr') dr'
h q 0 The beam emerging (after scattering) along
the incidence direction (θ = 0) contains fewer
dσ particles then the incident beam, since a number
So = |f (θ)|2
dΩ of particles have scattering in various directions.
∞ This decrease in the number of particles is
∫
4µ2
=
h4q2
| 0
r' V (r') sin (qr') dr'|2 measured by the total cross section σ, that is, the
number of particles removed from the incident
Validity of first Born approximation— beam along the incident direction is proportional
to σ, or equivalently to the imaginary part of f (0).
µ ∞

h2 k
| ∫0 V (r') (e 2ikr' – 1) dr'| << 1 Scattering of Identical Particles
Partial wave analysis—We assume here the Classically, the cross section for the
potential to be spherically symmetric. The scattering of two identical particles whose
angular momentum of the incident particle will interaction potential is central is given by
therefore be conserved, a particle scattering from σel (θ) = σ (θ) + σ (π – θ)
a central potential will have the same angular The differential cross section is
momentum before and after collision. dσ
δl is a real angle which vanishes for all values = |f (θ) + f (π – θ)2 = |f (θ)| 2
dΩ boson
of l in the absence of the scattering potential (i.e. = |f (θ)|2 + |f (π – θ)|2 + f (θ) *
V = 0), δl is called the phase-shift. f (π – θ) + f (θ) f* (π – θ)
∞
f (θ) = ∑ fl (θ) = |f (θ)|2 + |f (π – θ)|2
l=0 + 2 Re [f * (θ) f (π – θ)] …(3)
∞
1
=
2ik ∑ (2l + 1) Pl (cos θ) (e2iδl – 1) In sharp contrast to its classical counter part,
l=0 eqn. (3) contains an interference term 2 Re [f * (θ)
∞ π
=
1
∑ (2l + 1) eiδl sin δl Pl (cos θ) …(2) f(π – θ)]. Note that when θ = , we have
k l=0
2
π 2
Where fl (θ) is known as the partial wave
amplitude. From equation (2) we can obtain the
( ) | ( )|
dσ
dΩ boson
=4 f
2
, this is twice as large as
the classical expression (which as no interference
differential and the total cross section.
π
dσ 1 ∞ ∞
= |f (θ)| 2 = 2 ∑ ∑ (2l + 1) (2l' + 1)
( ) [ ( )]
term);
dσ
dΩ Cl
=2 f
2
2
.
dΩ k l = 0 l′ = 0 If the particles were distinguishable, the
ei (δl – δl') sin δl sin δl' × Pl (cos θ) Pl ′ (cos θ) differential cross-section will be four times
π 2
σ = ∫dΩ
dσ
dΩ
smaller, ( ) | ( )|
dσ
dΩ dist
= f
2
.
1
∞ ∞ The wave of two spin particle is known to
4π
σ= ∑ σl = k2 ∑ (2l + 1) sin2 δl 2
l=0 l=0 be either symmetric or anti symmetric when the
 Physical Sciences | 307

spatial wave function is symmetric, that is the two ^ |φ >+

particles are in a spin singlet state, the differential
So E n = E0 + < φn |H p n ∑
m≠n
cross-section is given by ^
|< φm| Hp | (φ n ) >|2
dσs + ……
= |f (θ) + f (π – θ)|2 En (0) – Em(0)
dΩ
Validity of the time-independent perturba-
But when two particles are in spin triplet tion theory—For perturbation theory to work,
states, the spatial wave function is anti symmetric,
the corrections it produces must be small;
hence convergence must be achieved with the first two
dσa corrections :
= |f (θ) – f (π – θ)|2
dΩ The expansion parameter must be small :
If the incident particles are unpolarized, the
various spin states will be equally likely, so the
< φ | H ^ |φ 

triplet state will be three times as likely as the E (0)m – E0 (0)n < < 1 (n ≠ m)
n m
singlet : If the unperturbed energy levels En (0) and
3 dσa 1 dσs (0)
dσ
dΩ | = +
4 dΩ 4 dΩ
fermion
E m were equal (i.e. degenerate) then condition
would break down. Degenerate energy levels
3 1 requires an approach that is different from the non
= |f (θ) – f (π – θ)2 + |f (θ) degenerate treatment.
4 4
+ f (π – θ)|2 Fine structure and the anomalous Zeeman
= |f (θ)|2 + |f (π – θ)|2 effect spin-orbit coupling—
– Re |fx (θ) f (π – θ)| ^ = 1 1 dv^ – –
Hso · S· L
π π 2 2
2m ec r dr
where θ = , we have
2
dσ
( ) [ ( )]
dΩ fermion
= f
2
2
;
Using perturbation theory to calculate the
this quantum differential cross section is half the contribution of the spin-orbit interaction in a
∂σ π
classical expression,
∂Ω cl( ) [ ( )]
=2 f
2
2
, and hydrogen atom.
→
four times smaller than the expression ^ = p 2 – e2 +
H
e2 ^S· ^L
corresponding to the scattering of two identical 2me r 2m ec2r32

π
bosons, ( )
dσ
dΩ boson
=4 f
2 [ ( )]
2
. = H^ +H ^
0 so
The energy levels of hydrogen Eso(1) is

 3
Approximation Methods for Stationary
States |En (1)| α2  j (j + 1) – l (l + 1) – 4
Es0(1)=  
Time-independent perturbation theory— n l (l + 1) (2l + 1)
First order correction— h e ~ 1
α = = –
^ φ > meca0 hc 137
E (1) = <φ |W|
n n n
h –e2
So ^ |φ >
En = En (0) + < φn |H Since a0 = and hence En(0) =
p n mee 2 2a0 n2
and first order wave function is α mec
2 2
^ |φ > = – . we can express
< φn |H 2n2
|ψn > = |φn > + ∑ p n
1 φn >
n ≠ 0 En
(0) – Em(0)  3 
α 4m c2  j (j + 1) – l (l + 1) – 
4
Second Order Correction— Eso(1) =
e
 
^ ψ (1) > 2n3  l (l + 1) (2l + 1)
E (2) = <φ |W|
n n n
Relativistic correction—
^ |φ > | 2
|<φ | W α2z2
En (2) = ∑ E m(0) – E n(0)
m≠n n m
[
En =z2 E n (0) 1 +
n ( 2
–
2l + 1 4n
3
)]
308 | Physical Sciences

The Variational Method i 

exp  ∫ p (x') dx' 
c'2
ψ2 WKB(x) =
 h 

√
In variational method, one does not attempt to x
p (x)
solve the eigen value problem.
 
+ 2 exp  – ∫ p(x') dx' 
c" i
^ |Ψ> = E|ψ>
H
 h 
Then we find

√ p (x)
x1 < x < x2
δ E (ψ) = 0
 1 x 
exp  – h ∫x |p (x')| dx',
c3
where E (ψ ) is the expectation value of ψ3 WKB =
energy in the state | ψ > 
√ |p (x)|  2 
^
<ψ|H|ψ> x > x2
E (ψ) = The connection of ψ3WKB (x) to ψ2WKB (x)
<ψ|ψ>
If | ψ > depends on a parameter α, E (ψ) also  1 x2 p (x') dx' + π 
∫
c2
sin 
4 
depends on α . The minimum value of E (ψ) ψ2WKB =
provides on upper limit approximation for the true 
√ p (x) h x
energy of the system. Connection of ψ1WKB (x) to ψ2WKB (x)
1 π
∫x1 p (x') dx' + 4 
Variational method is particularly used for D x
ψ2WKB = sin  h
determining the ground state energy and its eigen

√ p (x) 
state without explicitly solving the Schroedinges
equation. Quantization of the energy level of the
Bound states—
Note that for any (arbitrary) trial function | ψ
∫x1 √
x2
> we choose, the energy E as given by above ƒ p (x1 En) dx) = 2
O 2m (En – V (x)) dx
equation is always larger than the exact energy E0 .

E =
< ψ |H| ψ >
≥ E0
= (n + 12)h
<ψ|ψ> Bound states for potential wells with one
The WKB (Wentzel – Kramers – Brillouin) rigid wall—
Methods—This method is useful for approximate x2
treatments of systems with slowly varying
potentials, that is potential which remain almost
∫ x1
p (x) dx = n + ( )
3
4
πh

constant over a region of order of the de-Broglie Remark—

wavelength. In the case of a classical systems, this We may state that the phase factor α of the
property is always satisfied since the wavelength WKB solution is in general equal to—
of a classical system approaches zero. The WKB (1) Zero for turning points located at the rigid
method can thus be viewed as a semi classical walls.
approx. π
(2) for turning points that are not located at
Validity of WKB approximation— 4
the rigid wall.
 –    h 
d λ (x) =  d 
Bound states for potential wells with two
 <<1
dx dx  p (x) rigid walls—
x2
– h
λ = and the above condition means that the
∫ x1 p (x) dx = nπh, n = 1, 2, 3
p
rate of change of de-Broglie wavelength is small. Relativistic Quantum Mechanics
Bound states for potential wells with Klein Gordan equation—
no rigid-walls—As x → –∞ and x → + ∞ –
[ 2 – k2] ψ (x, t) = 0
respectively
2 = ∇2 – 1 ∂
2
 1 α1 
exp  – h x |p (x')| dx',
where
∫
c1
ψ1 WKB(x) = c2 ∂t2

√ |p (x)|   mc
and k = h0
x < x1
 Physical Sciences | 309

Its Lorentz covariance can be readily be seen Multiply by β from the right
by introducing momentum four-vector with its αxββ = –βαxβ
components
{pµ} = {p1 , p2 , p3 , iE/c} αx = –βαxβ
∂ Trace αx = –Trace (βαxβ)
where pµ = – ih
∂xµ [Since trace (ABC) = Tr (CAB) = Tr (BCA)]
where {xµ} = (x1 = x; x2 = y; x3 = z; x4 = ict) So Tr (α x) = –Tr (α xββ) = –Tr (α x)
Then = ∂ µ dµ 2 Trace (αx) = 0
∂
with ∂µ = , so, K. G. equation can be or Tr αx = 0
∂xµ
written as 4. We know that the trace is just the sum of
  m0c 2  all eigen values, it follows that +1 and –1 eigen
 ∂µ∂µ ·  h   ψ = 0 values must occur the same number of times. This
shows that each of the matrices must of even
where, energy is given by dimension n.
E = – c √
p2 + m20 c2 Probability and current density—Probabi-
lity density
KG equation fell into disgrace if it is consi-
–
dered as state equation but if we consider rather p = ψ+ψ
as the field equation of the quanta associated
with spinless particles. and probability current density
The direc equation—In order to avoid the S = Cψ+ αψ
difficulties faced by Klein-Gardan equation as the  ψ1 
result of non-definite probability density, we  ψ2 
p = (ψ1 * , ψ2 * , ψ3 * , ψ4 * )  ψ 
–
follow the historical path adopted by dirac in So
seeking a relativistic covariant equation linear in  3 
ψ4
time and spatial derivatives. This approach leads
to a four component wave equation which a = ψ1 * ψ1 , ψ2 * ψ2 , ψ3 * ψ3 , ψ4 * ψ4
1
accounts for many properties of spin – particles. = ∑ |ψi|2
2 i

^ has been defined as which is positive definite and hence can be used
Hamiltonian operator H as the probability density.
^ = c (–α·–p) + ^βm c2
H –
0
S → probability current density specifying
where ^α and β^ are known as Direc matrices. the mode of the probability flow. The operator cα
^ can be interpreted as velocity.
Properties of α
^ and β
Matrices Free particle solutions—
1. Since αx2 = αy2 = αz2 = β2 = 1
∴, The eigen values of matrices must be ± 1,
 10 
2. Since all these matrices anti commute (no  cp 
E + mc 
z
UI = 2
two matrices commute), therefore only one of the
four matrices, can be diagonal. Here we take β to
be diagonal, so others elements are non-diagonal
 cp 
+
E + mc2
 01 
matrices.
3. Trace α i = 0, i = x, y, z
Trace β = 0  cp 
E + mc 
–
For trace α x, consider UII = 2
αxβ + βαx = 0  cp 
z
αxβ = –βαx E + mc2
310 | Physical Sciences

 cp 
z
2
electron. The last two terms are of order
h
c
and
 cp 
E + mc
hence neglected in non-relativistic case.
E + mc 
+
UIII = 2 Electron spin—Let us analyse the laws of
 1 
 0
conservation of angular momentum in a central
force field, V (r).

 cp 
In non-relativistic theory the orbital angular
– momentum, commute with non relativistic
2
 cp 
E + mc Hamiltonian.
E + mc 
z [L, H] = 0
UIV = 2

 0  However, in relativistic quantum mechanical

angular momentum operator does not commute
1
with dirac Hamiltonian even for a free particle. So
where E = –mc 2 (negative energy) the orbital angular momentum is not a constant of
motion.
Ep = mc 2 (positive energy solution)
dL –
ih x = [Lx1 H] = [Lx, cα·p + βmc2 + V (r)]
Dirac Equation in Electromagnetic dt
Field (Megnetic Moment) = ihc {αyp2 – αzpy]
dLx
The Direc equation in electromagnetic field is ≠ 0
dt
given by
or Lx ≠ constant
– –
[(E – eφ)2 – (cp – eA)2 – m2c4 + ihecσ'. H – This means that we must find another ope-
–– rator such that the commutator of its x-component
ichc α·E] ψ = 0. with H is the negative of right hand side, the sum
The first three terms are precisely the same as of this operator and L is then a constant of motion
in relativistic wave equation in electromagnetic and can be interpreted as the angular momentum.
field. Here spin angular momentum S is defined as
Physical significance of last two term—The h
eh S = ∑x
2
term σ' represents magnetic dipole moment Spin-orbit coupling—The spin orbit coupl-
2mc
of electron. ing energy is
ieh → 1 1 dV – –
and α → imaginary electric moment of an Us – o =
2m2c2 r dr
S· L
2mc

Solved Questions
Q. 1. The wave function of a free particle is Applying the condition of normalization
– x2
represented by Ψ(x ) = N exp
2a2( + ikx . ) +∞
∫ –∞ Ψ* Ψ dx = 1
Find out the value of constant N and the linear +∞
probability density. ⇒∫ –∞ (Ne–x2/2a2 e–ikx) (Ne–x2/2a2 eikx)dx = 1
+∞
Solution : Given ⇒ N2 ∫ –∞ e–x2 /a2 dx = 1
Ψ(x) = Ne –x2 /2a2 eikx Using the definite integral

∫ –∞ e–αx2 dx=√
α
Its complex conjugate be +∞ π
Ψ* (x) = N e–x2 /2a2 e–ikx

THERMODYNAMIC AND STATISTICAL
PHYSICS
Zeroth Law of Thermodynamics
There exist a scalar quantity called tempera-
ture which is a property of all thermodynamic
systems such that the equality of temperature is
necessary and sufficient condition for thermal
equilibrium.
First Law of Thermodynamics
As a result of its interaction with the environ-
ment, the thermodynamic system can receive or
give up a certain quantity of heat d Q, can do work
or work can be done on it, dW. In all cases a
change in the internal energy dU must be equal to
the difference between the amount of heat. dQ, the
system gains from outside and work dW = PdV it
does in opposing the external pressure P.
dU = d Q – d W
= d Q – PdV
Here dV is the change in the volume of
system. However this law does tell us about the
quantity of heat change into work and the direc-
tion in which the changes take place.
Second Law of Thermodynamics
The amount of heat obtained by the system in
a reversible process causes an increase in its
entropy by d S = dQ/T, where T is the temperature
at which the system takes up heat.
Third Law of Thermodynamics
As T → 0, the entropy of a system S (T, a3 )
causes to depend on a3, i.e., S is a fixed quantity
for all substances. The parameter a 3 are the
generalized coordinates of external bodies viewed
as additional variables of the system.
Thus, the method of thermodynamics rests on
the above mentioned laws, using the fact that S
and U are state functions of the system for which
Unit-5
∫
O dU = 0
∫ dS
O = 0
i.e., d U and d S are total or exact differentials.
The differential d Z = X d x + Ydy is an exact
differential it
∂X ∂Y
( ) ( )
∂y x
=
∂x y
Entropy
Entropy S is a state function defined by the
equation
dS =
1
T ∫dQ
where d Q is the increment in the heat energy
supplied to the system and T the absolute tem-
perature.
The entropy change d S in an isolated system
is zero for a reversible process and positive for all
irreversible processes.
The first law of thermodynamics confirms the
fact that work may be dissipated completely into
heat whereas heat may not be converted in work
totally. This shows the one sideness of nature. The
second law is law of spontaneous change. Heat
flows from hotter to colder region. Gases seep
from high pressure to low pressure region. The
deduction can be made in terms of time. Does
time move forward ? Because the entropy always
icreases, the physical universe is moving from a
less probable to a more probable state, times
arrow points forward. Another question arises if
entropy is increasing the difference in temperature
will vanish or there will be an uniform tempera-
ture. In the oscillating model of universe it is
supposed that entropy inequality is reversed at the
universe’s greatest extent so that during contrac-
tion entropy tends to a minimum.
 Physical Sciences | 371

Also acquisition of heat by a substance 2. Isobaric process—These are the process in

implies more energy for its atoms and molecules. which the pressure remains constant during entire
Entropy is linked to the atomic constitution of process.
matter by the equation, Work done during isobaric process.
S = k log W W = P(V2 – V1)
where W is the probability of distribution of the ∆Q = C P∆T
particles of matter among various energy states = C V∆T + P∆V
and k is Boltzmann’s constant.
3. Isochoric process—In these process
Combination of First and Second Law volume remains unchanged during the process.
of Thermodynamics Work done in isochoric process will be
The definition of entropy as an outcome of
second law of thermodynamics can be given by
W = ∫PdV = 0 (Q dV = 0)

δQrev so the first law gives

dS = …(1) ∆Q = ∆U
T
Also the first law of thermodynamics for for one gm of ideal gas at constant volume
reversible conditions can be given by ∆Q = C V∆T
δQrev = dE + PdV …(2) where C V = Specific heat of the gas at
where PdV represents the work wasted against constant volume
confining pressure on combining equations (1) ∆T = Change in temperature
and (2) we can obtain ∴ ∆U = C V∆T = ∆Q
TdS = dE + PdV …(3)
4. Adiabatic process—The process in which
Equation (3) is known as combined form of no exchange of heat takes place between the
first and second law of thermodynamics. system and surroundings. These kind of process
Ideal Gas Processes takes place quickly or suddenly. The curve
showing the variation of volume with pressure
1. Isothermal process—If the process is under adiabatic condition is called an adiabatic.
carried out in such a way that the temperature
remains constant throughout, the process is said to For an adiabatic change
be isothermal e.g., melting and boiling. ∆Q = 0
PV = RV = constant Hence ∆U = ∆W = 0
(for one mole of ideal gas) or ∆U = –∆W
and PV = nRT i.e. change in internal energy is mechanical
(for n mole of an ideal gas) work done or we can say that the work, in an
For isothermal process, adiabatic process, is done at the expense of inter-
since temperature remains same ⇒ ∆U = 0 nal energy.
Thus, first law becomes The adiabatic equation of a perfect gas is
∆Q = ∆W given by
Total work done during an isothermal PVγ = constant (Poisson’s law)
process. where P and V denote pressure and volume
respectively and γ is the ratio of specific heats at
∫
V2
W = PdV constant pressure and constant volume respec-
V1
tively.
∫
V2
RT
= dV C
VV1 γ = P.
V CV
= RT ln 2 If we take T as temperature, then the above
V1
expression will take the form
 V2 TVγ – 1 = constant
=  2·303 log10 V  RT (using PV = RT)
 1 and TP (1 – γ)/γ = constant
372 | Physical Sciences

Work done in adiabatic process, Maxwell’s Thermodynamical Relations

1 A thermodynamic state of a system can be
W = [P V – P 1 V1]
1–γ 2 2 described by thermodynamic coordinates such as
1 P, V, T, U, Q, S etc. Out of which pressure P,
= [RT 2 – RT1 ]
1–γ temperature T, volume V and entropy S are basic
R independent coordinates while other are combina-
W = [T – T2 ] tions of these. These coordinates are dependent of
γ–1 1
each other. By using the laws of thermodynamics
for n molecule of an ideal gas
Maxwell deduce six fundamental thermodynamic
nR
W = [T – T2 ] relations which are applicable to all thermo-
γ–1 1 dynamic systems.
The state of a system may be defined as any
To find the number of molecules is a given pair of quantities viz., P, V, T and S. We can
mass of Gas—If m gm is the given mass of a choose the most suitable pair for the solution of
gas and M gm is the molecular mass of the gas thermodynamic problem. The quantities appeared
m in the pair all independent variables from first law,
then number of moles in m gm =
M we have
e.g. for 1 gm of hydrogen d Q = dU + dW
Mass of gas or d Q = dU + PdV
Number of moles =
Molecular mass or d U = d Q – PdV …(A)
1 From second law, we have
= = 0·5
2 d Q = TdS …(1)
and the combined relation
General relation between pressure and volume
of a gas— d U = TdS – PdV …(2)
PVn = constant Suppose S, U and V are functions of two
where n is a constant depending upon the process independent variables x and y.
(i) n = 0; the process is isobaric
(ii) n = 1; the process is isothermal
dS =
∂S
( ) ( )
∂x y
dx +
∂S
∂y x
dy 
(Boyle’s law)
dU = (∂U∂x ) dx + (∂U∂y ) dy 

(iii) n = CP/CV = γ; the process is adiabatic y x
…(1A)

dV = (∂V∂x ) dx + (∂V∂y ) dy
y x

n=0 Isobaric The subscripts x and y denote the constant
coordinates. From equation (1A) and (2) we get
P
n=1 Isothermal (∂U∂x ) dx + (∂U∂y ) dy =T (∂S∂x) dx + (∂S∂y) dy
y x y x

n=γ  ∂V ∂V 
– P  ( ) dx + ( ) dy  …(1B)
Adiabatic
V
 ∂x y ∂y  x
[Graphical comparison of isobaric, isotherm
We can compare the coefficients of dx and dy
and adiabatic] on both the sides
From the curve;
∂U ∂S ∂V
Slope of adiabatic
Slope of isotherm
= γ
( )
∂x y
= T
∂x y ( ) ( )
–P
∂x y
…(3)

Obviously, adiabatic curve is steeper than

isotherm.
and (∂V∂y ) x
∂S
∂y
∂V
= T( ) –P( )
∂yx x
…(4)
 Physical Sciences | 373

On differentiating equation (3) w.r.t. y and On substituting in equation (5)

equation (4) w.r.t. x we obtain
∂2 U ∂T ∂S ∂2 S (∂S∂V) = (∂P∂t )
∂y∂x
= ( )( ) ∂y x ∂x y
+T
∂y∂x
T

1
V

˙.˙ dS = dQ
∂P ∂V ∂ V 2 T
–( ) ( ) –P
∂y ∂x ∂y∂x

∂ U 2 ∂T ∂S ∂ S
x y

2
(∂Q
∂V) T
= T (∂P∂T) V
and = ( ) ( ) +T
∂x∂y ∂x ∂y ∂x∂y y x Second relation—Taking T and P as inde-
pendent variables
∂P ∂V ∂ V 2
–( ) ( ) –P
∂x ∂y ∂x∂y y x
x = T, y = P

The change in the internal energy d U brought ∂T ∂P

Then = 1, =0
about by changing V and T, whether V is changed ∂x ∂y
by d T first and T is changed by dV later or vice- ∂T ∂P
versa remains the same, which shows that d U is a = 1, =0
∂y ∂x
perfect differential, i.e.,
On substituting in equation (5)
∂2 U ∂2 V
∂y∂x
=
∂x∂y (∂S∂P) T
= – (∂V
∂T ) P

(∂T∂y) (∂S∂x) + T ∂x∂y

∂ S ∂P
∂y
∂V
2
–( ) ( ) –P
∂x
∂ V
∂y∂x
2

∂T
y

∂S ∂ S
x y

2
or (∂Q
∂P ) T
∂V
= – T( )
∂T P
= ( ) ( ) +T
∂x ∂y ∂x∂y y x Third relation—Taking S and V are indepen-
dent variables
∂P ∂V ∂ V 2
–( ) ( ) –P x = S, y = V
∂x ∂y ∂x∂y y x
∂S ∂V
∂2 S ∂2 S ∂2 V ∂2 V Then = 1, = 1
˙.˙ = ; = ∂x ∂y
∂y∂x ∂x∂y ∂y∂x ∂x∂y
∂S ∂V
and = 0, = 0
∂T ∂S ∂P ∂V ∂y ∂x
∴ ( ) ( ) –( ) ( )
∂y x ∂x y ∂y x ∂x y Substituting in equation (5)
∂T ∂S
∂y
∂P
∂x
y
∂V
= ( ) ( ) – ( ) ( ) …(5)
∂x ∂y x y x
(∂T∂V) S
= – (∂P∂S) V

where x and y can be any two variables out of P,

V, T and S. or (∂T∂V) S
= – T( )
∂P
∂Q V
First relation—Taking T and V as inde-
Fourth relation—Taking S and P as inde-
pendent variables pendent variables
x = T, y = V x = S, y = P
∂T ∂V ∂S ∂P
Then = 1, =1 Then = 1, =1
∂x ∂y ∂x ∂y
∂T ∂V ∂S ∂P
and = 0, =0 and = 0, =0
∂y ∂x ∂y ∂x
374 | Physical Sciences

Substituting in equation (5)

where (∂Q
∂V)
is the quantity of heat absorbed per

(∂T∂P) = (∂V
S∂S ) P
T
unit increase in the volume at constant tempera-
ture. This quantity is termed as latent heat. If ∂Q is
the latent heat, L for the change in volume of
or (∂T∂P) = T (∂V
S ∂Q) P
substance from V1 to V 2 at constant temperature
then
Fifth relation—Taking P and V as inde- ∂Q = L
pendent variables
and ∂V = V2 – V1
x = P, y = V Putting in equation (1) we can obtain
∂P ∂V
Then
∂x
= 1,
∂y
=1 L
V2 – V1
= T (∂P∂T) V
∂P ∂V L dP
and = 0, =0
∂y ∂x or
V2 – V1
= T
dT
Substituting in equation (5)
dP L
or = …(2)
∂T ∂S ∂T ∂V dT T (V2 – V1)
( )( ) ( )( )
∂P V ∂V P
–
∂V P ∂P V
=1
Equation (2) is referred as Clapeyron’s latent
heat equation.
Sixth relation—Taking T and S as inde-
pendent variables (a) Effect of change of pressure on the
melting point—When a solid is going through
x = T, y = S melting there is a change in its volume.
∂T ∂S dP
Then = 1, =1 If V2 > V1, is a positive quantity so the
∂x ∂y dT
∂T ∂S melting point of the substance will increase with
and = 0, =0 the increase in pressure.
∂y ∂x
dP
Substituting in equation (5) If V2 < V1 , is a negative quantity (for ice
dT
etc.) so the melting point of such substances will
(∂P∂T) (∂V
S
∂P ∂V
∂S ) (∂S ) (∂T )
–
T
=1
T S
decrease with increase in pressure and vice-versa.
(b) Effect of change of pressure on boiling
These are Maxwell’s six thermodynamic point—When a liquid is boiled there is always an
equations. Sometimes known as thermodynamic dP
increase in its volume, i.e., V2 > V1 thus is a
tools also. dT
positive quantity always. So with increasing
Clausius-Clapeyron’s Latent Heat pressure, the boiling point of a substance increases
Equation and vice-versa. The liquid will boil at a lower
temperature under reduced pressure.
It is important equation to tell us the effect of
pressure on melting and boiling points of sub- Thermodynamic Potentials
stances.
Internal energy—Due to the thermal motion
From Maxwell’s first thermodynamic relation of the particles and mutual attractive forces
between them every system possesses some
(∂Q
∂V) T
= T (∂P∂T) V
…(1) intrinsic energy (K.E. + P.E. molecules). This
energy is called internal energy (U).

From first law of thermodynamics,
dQ = dU + dW
or dU = dQ – dW
In terms of entropy (S)
dU = TdS – PdV (Q ∂Q = T∂S)
The internal energy is an unique function i.e.
the change in internal energy dU is independent of
the process chosen but depends on initial and final
states only.
Enthalpy—The absolute value of internal
energy U cannot be calculated from classical
thermodynamics. Therefore, a new thermo-
dynamic function, called enthalpy (H) introduced.
Enthalpy and internal energy are related by
following relations—
H = U + PV
E, P and V are the properties of the state of
the system, it is evident that H is a function of
state of the system and its value is independent of
the path by which this change has been
accomplished.
Free Energy Work Functions
In most of the thermodynamic process,
neither the energy nor the entropy are held
constant. It becomes, therefore, essential to
introduce two new functions involving energy and
entropy. These are—
(a) The Helmholtz free energy or work
functions F—The work function A can be defined
as maximum work content of a system and is
given by
F = U – TS
Here U and S denote the internal energy and
entropy of a system respectively. Therefore, F
must be a state function and d F will be a perfect
differential.
(b) The Gibbs free energy work function
G—The second function also called Gibbs free
energy is defined by the equation
G = H – TS (where H is enthalpy)
or G = (U + PV) – TS
= (U – TS) + PV
G = F + PV
Physical Sciences | 375
This equation gives an alternative definition
of free energy and also relates work function F
with free energy G. As F, H and S are thermo-
dynamic properties of the system, it follows that G
must be a state function like work function F.
Subsequently d G like d F would be a perfect
differential.
Importance of Thermodynamic Poten-
tials or Functions
Thermodynamic potentials are of practical
importance in studying the equilibrium conditions
of a thermodynamic system, e.g., a system of
which the temperature and volume are kept con-
stant, is in equilibrium when the free energy F
is a maximum. It is because
d F = – PdV – Sd T
In an isothermal-isochoric process d V = 0,
dT = 0 and hence,
dF = 0
Thus, F = F min
When T = Constant, V = Constant
So a system at constant temperature and
volume is in stable equilibrium when Helmholtz
free energy F (= U – TS) is a minimum. It is
analogous to a mechanical system which is in
stable equilibrium when the potential is minimum.
Similarly a system at constant pressure and
temperature is in stable equilibrium when the
Gibbs’ potential G (= H – TS) is the minimum. It
is because
d G = Vd P – Sd T
In an isothermal isobaric process
dT = 0
dP = 0
Hence dG = 0
Thus, G = Gmin
When T = Constant, P = Constant
Gibbs-Helmholtz Equation
The Gibbs’ free energy function G is defined
by the relation
G = H – TS …(1)
On differentiating it gives
d G = d H – TdS – Sd T …(2)
376 | Physical Sciences

Since H = E + PV which on differentiation It is clear that above the temperature 31·1°C

gives the CO2 cannot be liquified. This temperature is
d H = d E + PdV + Vd P …(3) known as critical temperature.
“The temperature below which the gas can be
Eliminating d H between (2) and (3), we get liquified by the application of pressure alone is
d G = d E + PdV + Vd P – TdS – Sd T called the critical temperature.” Above this tem-
perature the gas cannot be liquified howsoever
Also TdS = d Q
high the pressure may be. The coordinates corre-
⇒ TdS = d E + PdV sponding to P are called critical temperature (TC ),
Therefore, critical pressure (PC ) and critical volume (VC).

d G = d E + PdV + Vd P – d E – PdV – Sd T At critical point, the liquid and vapour states

are indistinguishable. Their densities are equal and
⇒ d G = Vd P – Sd T …(4) their mixture is homogeneous with high com-
at constant pressure, pressibility. The surface tension vanishes at
critical point.
dP = 0
This point P can be obtained by Vander
Equation (4) becomes Waals equation,
d G = – Sd T RT a
P = –
∂G V – b V2
or ( ) ∂T P
= –S
at critical point
d2 P
= 0
dt2
We can introduce the value of S in equation
(1), i.e., which gives VC = 3b

G = H+T (∂G
∂T ) P
PC =
a
27b2
This is known as Gibbs-Helmoltz equation. 8a
TC =
27Rb
Critical Constants and the relation between Boyle temperature (TB ),
In terms of a and b. temperature of inversion (Ti) and critical tempera-
The isotherms of CO2 at various temperatures ture be
as obtained by Andrew are shown below : 2a
Ti = ,
Rb
a
TB =
Rb
TB is the temperature at which the gas follows
Boyle’s law, i.e.,
d (PV)
= 0
dP
From above
Ti > TB > TC
Law of Corresponding States
The portion AB represents the gaseous state Two gases are said to be in corresponding
of CO2 and between A and B gas obeys Boyle’s states if they have the same values of the ratio of
law. From B to C, CO2 changes from gaseous to their actual pressure, volume and temperature to
liquid state and CD represent the liquid form of the critical pressure, critical volume and critical
CO2 . temperature.
 Physical Sciences | 377

P1 P or U1 + P1 V1 = U2 + P2 V2
i.e. = 2
P C1 P C2 ⇒ (U 1 – U2) = P 2 V2 – P 1 V1
V1 V2
= Vander Waals Gas in Joule-Kelvin
VC 1 VC 2 Effect
T1 T2 The fall in temperature is given by—
and =
TC1 TC2 (P1 – P 2 ) 2a

Let
P
= α
dT =
CP (
RT
–b )
PC Let us define Joule-Kelvin coefficients as
V follows :
= β
VC
T
µ = (dTdP) H
=
1
CP [RT2a – b]
and = γ
TC µ denotes the slope of iso-enthalpy curves on
Then the Vander Waals gas equation takes the T-P plane.
form a, b, c and d are the maxima of their respec-
tive curves. The dotted lines shows locus of
(α + β3 ) (3β – 1) = 8γ
2 inversion points and is called as ‘inversion curve’.

This is known as reduced equation of state.

When two gases have the same values for any two
of the three reduced variables α, β and γ they also
have the same values for the third variable and
said to be in corresponding states.

Joule-Kelvin Effect or Joule-Thomson

Effect
Porous plug consist of compressed porous
material like cotton, wool or silk fibre and thus Now, since P2 < P1; dT and also µ will be—
equivalent to number of orifices in parallel.
2a
Throttle molecules gets further apart from one (i) Positive if –b > 0
RT
another under a pressure drop. 2a
⇒T< i.e. cooling will take place.
In the experiment the gas is forced at constant Rb
pressure through a porous plug thermally insulated 2a
from outside. The gas on passing through the plug (ii) Zero if –b = 0
RT
under a large drop of pressure gets throttled. Thus, 2a
the gas has to do some internal work against inter- Ti =
Rb
molecular attractions, if these exists. Since no heat
is supplied, this work done at the cost of internal This temperature is called temperature of
energy. Hence, the gas passing through the plug inversion and is denoted by Ti. Thus, the tempe-
will suffer a fall in temperature and so the cooling rature of the gas remain unchanged if T = Ti.
was demonstrated. 2a
(iii) Negative if –b < 0
In this experiment, we find that the initial heat RT
constant is same as the final one, i.e. it is iso- 2a
enthalpy effect. or T >
Rb
H1 = H2 T > Ti.
378 | Physical Sciences

Hence, the gas will be heated when it is Second order—Involves no transference of

forced through the porous plug if the initial tem- heat and there is no change in volume. Thus, it
perature is greater than temperature of inversion. takes place with no change in entropy and volume
at constant temperature and pressure e.g.
This explains why hydrogen and helium,
whose temperature of inversion are much below (i) Transition of liquid helium I to liquid
the room temperature, show heating effect at helium II.
ordinary temperature.
(ii) Transition of ferromagnetic to para-
This effect is different from adiabatic expan- magnetic material at Curie point—
sion in the terms that in adiabatic expansion, only
cooling takes place because the work is done at
the expense of internal energy.

Phase Transition
First order—The one in which Gibb’s
function with respect to pressure and temperature
change discontinuously at the transition point.
However, the value of the Gibbs’ function is the
same in both the phases at equilibrium. The
Clausius-Clapeyron latent heat equation holds
good for first order transition.

(Entropy)
δG
—–
δT

Law of Equipartition of Energy

It states that the total kinetic energy of a
(Volume) dynamical system consisting of a large number of
particles in thermal equilibrium is equally divided
δG
among its all the degrees of freedom and energy
—– 1
δP associated with each degree of freedom is RT,
2
where k is Boltzmann’s constant and T is absolute
P temperature of the system.
 Physical Sciences | 379

Maxwell’s Distribution Law of and corresponding maximum probability

Velocities
Basic assumptions—(i) The gas consists of
P max = [2πkT
m
] 1/2

molecules with all possible velocities between –∞ Thus Pmax increases as m increases and T
to +∞. decreases
(ii) In steady state, the density remains
uniform.
(iii) Though the velocities are continuously T1
changing due to collisions, the number of
molecules having velocities in a definite range
remains uniform. T2 T1 < T2 < T3
Let m → mass of one molecule
vx → Velocity component in X-direction T3
vy → Velocity component in Y-direction.
vz → Velocity component in Z-direction (– ve)vx vx (+ ve)
k → Boltzmann’s constant
The area under each curve is unity.
and T → Absolute temperature.
Then number of molecules having velocity Maxwell’s Distribution Law of Speed
coordinates in the range vx to v x + dv x, vy to v y + dv y
The number of molecules having speed v
and vz to vz + dvz is given by
to v + dv is given by
ni dvx dvy dvz = N ( )m 3/2
2πkT
–m
Nvdv = 4πN ( ) m 3/2 2 –mv2/2kT
2πkT
ve dv
(v2x + v2y + v2z)
e 2k T dvx dvy dvz and the probability is
and the number of molecules having velocity com- Nvdv
Pvdv =
ponents in the range vx to vx + dvx is N
mv 2x
ni dvx = N ( )
m
2πkT
–
e 2kT dvx. = 4π ( ) m 3/2 2
2πkT
v × e–mv2/2kT dv

The probability that a molecule will have x Most probable speed—

component of velocity in the range vx to vx + dvx is dP v
given by = 0, which gives;
dv
n dv

√
P(vx)dvx = i x 2kT
N vmp =
mv 2 m
= ( )
m 1/2 – kTx
2πkT
e dvx. Average speed—

∫
– ∞
Average Velocity v = 0 v Pv dv

√8kT

It is clear that P(v x) is symmetrically dis- –
tributed about the value v x = 0 v =
mπ
∫
– +∞
v x = – ∞ vx P(vx)dvx Root mean square speed—

∫
–2 ∞
Maximum Probability v = 0 v2 Pv dv
For probability to be maximum; 3kT
=
dP(vx) m
= 0
√3kTm

dvx
which gives vx = 0 or vrms =
380 | Physical Sciences

– But since all molecules are moving in all

Clearly vrms > v > vmp possible directions, therefore,
1
λ =
 2 πσ2n
√
Note—For transmission of sound through the
gas λ should be less than the wavelength of sound.
m kT/ρ
Now, λ = =
2
√ πσ2ρ  2πσ 2
√
At a given temperature where ρ is the density of gas.
1
Hence, λ ∝
ρ
λ ∝ T
1
λ ∝
P
1
λ ∝
n
Deductions— the probability that a molecule would traverse
(i) Fraction of molecules having speed greater distance x without colliding (or will collide after
than a particular speed increases as temperature traversing x is)
increases. P = e–x/λ.
(ii) vroot mean square > vaverage > v most probable
Black Body Radiations
(iii) More than half of the molecules have
speed greater than v mp . Provost suggested that everybody emits
radiations continuously and also receives from
(iv) Fraction of molecules having speeds surroundings at a rate depending on the nature of
within a given range of r.m.s. speed decreases as its surface and temperature.
temperature increases.
A body at higher temperature emits more
Mean Free Path radiations than it receives, consequently the tem-
perature falls. For the body at lower temperature it
The fact that molecules collides itself suggest happens the other way around.
that they are not point masses, rather they have a A black body absorbs all the heat radiations,
finite size. The path between two collisions is of whatever wavelength, incident on it and thus,
called free path. The average distance traversed by appears black.
a molecule between two successive collisions is
mean free path. Spectrum of black body radiations—

If n → Number of molecules per unit volume T1

v → Average speed of molecule T1 > T2 > T3 > T4
T2
and σ → Diameter of molecule.
Then, considering all other molecules at rest
the mean free path (λ) of molecule under consi- E
deration is T4 T3

1
λ = .
πσ 2 n λ
 Physical Sciences | 381

Observations— harmonic oscillators of molecular dimensions and

(i) Energy is not distributed uniformly in the
spectrum.
(ii) At a given temperature, the intensity
increases with wavelength and becomes maximum
at a particular wavelength and then decreases.
(iii) λm (wavelength for maximum intensity)
decreases with temperature increases.
(iv) The area under each curve increases with
temperature and is proportional to fourth power of
absolute temperature.
(v) Increase in temperature causes an increase
in energy emission for all wavelengths.
the exchange of energy between matter and radia-
tion does not take place continuously but
discontinuously and discretely as an integral
multiple of small unit of energy called photon
or quantum.”
Energy of each photon is
ε = hν where ν is the frequency and h is
Planck’s constant h → 6·62 × 10–34 J-sec.
Average energy of classical oscillator = kT
Average energy of a Planck’s oscillator
hν
ε= .

Wien’s Displacement Law

[ ( ) – 1]
exp
hν
kT
The number of modes of vibration per unit
Product of wavelength corresponding to volume within frequency range ν and ν + dν
maximum energy λ m and absolute temperature T
is constant i.e. uπν2
= dν
λmT = constant c3
= 0·2896 × 10–2 mK Also number of modes in wavelength range λ
Wien’s Distribution Law and λ + dλ

A 8π
uλdλ = f (λT)dλ = dλ
λ5 λ4
= AT5 f (λT)dλ
Rayleigh-Jeans Law
where A is a constant and f denotes a function.
Clearly energy density is proportional to fifth Energy density within frequency range ν and
power of temperature and inversely proportional ν + dν
to fifth power of wavelength. 8πν2 *
uν dν = dν kT
Wien’s law confirms to the experimental data c3
only for short wavelength range. 8πkT
or uλ d λ = dλ
Solar Constant λ4
Amount of solar energy received per minute This explains the experimental measurement
by unit area of a perfectly black surface placed at of the energy distribution at long wavelength
mean distance of the earth from the sun, in the fairly well; but it fails for short wavelengths.
absence of atmosphere and at right angles to the
direction of Sun’s rays. Statistical Thermodynamics
Typical value 2 cal/cm2/min. This is an approach to deal with applica-
or 1·38 kW/m2 tions of statistical mechanics to thermodynamics.
Statistical thermo-dynamics relate the microscopic
Quantum Theory of Radiations : properties of the individual molecules (moment
Planck’s Hypothesis of inertia, dipole moment etc.) with the macro-
scopic properties (molar heat capacity, polariza-
“A black body radiation chamber is filled up tion etc.) of a system having a large number of
not only with radiations; but also with simple molecules.
382 | Physical Sciences
It can be used for systems at equilibrium as
well as away from equilibrium.
Probability Theorems
(1) The number of ways in which N
distinguishable particles can be arranged in order
will be equal to
N!.
(2) The number of different ways in which n
particles can be selected from N distinguishable
particles irrespective of the order of selection will
be equal to
NC = N!
n
n ! (N – n) !
(3) The number of different ways in which n
indistinguishable particles can be arranged in g
distinguishable states with not more than one
particle in each state will be equal to
g!
n ! (g – n) !
Phase Space
In order to specify the state of the gas from
the molecular point of view, we shall know the
position and momentum of each of its molecules,
i.e., we must specify six quantities x, y, z, px, py, pz
for each of molecules. In three dimensional space,
we can easily consider an element of volume dx,
dy, dz. Similarly, we can plot momentum compo-
nents px, py and pz in three mutually perpendicular
directions and again consider three dimensional
volume dpx, dpy, dpz. Thus, we may imagine a six
dimensional space in which dx.dy.dz.dpx.dpy.dpz is
an element of volume and position of a point in
this space will be described by a set of six
coordinates x, y, z, px, py, pz. This six dimensional
space is known as phase space and element in this
space is called cell.
The meaning of a point in phase space can be
further understood by means of uncertainty
principle. This can be carried out by dividing a
phase into small six dimensional cells whose sides
are dx, dy, dz, dpx, dpy, dpz . Let us approach close
to the limit of a point in phase space, as the size of
cells is reduced. However volume of these cells
can be given by
d λ = dx.dy.dz.dpx.dpy.dpz
and according to uncertainty principle
dx.dpx ≥ h
dy.dpy ≥ h
dz.dpz ≥ h
Thus, we find
dλ ≥ h3
A point in phase space is actually considered
to be a cell whose minimum volume is of the order
of h3 . A particle in phase space may be understood
as being located somewhere inside the cell centred
at some location x, y, z, p x, py, pz instead of being
precisely at the point itself. The notation of a point
of infinitesimal size is surely acceptable either in
position space or in momentum space alone but
not in phase space because it will violate the
uncertainty principle.
Occupation Number
It is the number of systems in that particular
state. The set of occupation number is called a
distribution.
Microstates and Macrostates
A microstate of the system may be defined by
the specification of six coordinates x, y, z, p x, p y,
pz of each molecule of the system within the limits
of the dimensions of the cell in which its represen-
tative points lies.
On the other hand a macrostate of the system
may be defined by the specification of the number
of molecules or phase point in each cell of phase
space.
Many different microstates corresponds to the
same macrostate.
Let a particular microstate be specified by
stating that the phase points or molecules a, c, h
are in cell 1, b, e are in cell 2. f is in cell 3 and so
on as shown in the figure. The corresponding
macrostate is specified by merely giving the
number of phase points n1 (= 3) in cell 1, the
number n2 (= 2) in cell 2 and number n3 (= 1) in
cell 3 etc. If we interchange any two molecules
from different cells say a and f, we will have the
different microstates but same macrostate. If we
interchange the two molecules in the same cell say
a and c, we will have the same microstate as well
as same macrostate.
 Physical Sciences | 383

Cell 1 a, c, h n1 = 3 For general purposes, it is required to measure

energy level, relative to the ground state. It,
therefore, follows that equation becomes as—
∞
Cell 2 b, e n2 = 2 Q = Σ gi e–εi /kT
i=0

It is a dimensionless quantity. Its value

Cell 3 f n3 = 1 depends on the molecular weight, molecular
volume and temperature etc. Thus, the partition
function may be defined as the sum of the proba-
˙˙˙˙˙˙˙

˙˙˙˙˙˙˙

bility factors for different energy states or more

conveniently it can be stated as the way in which
Cell i dk ni = 2 the energy of a system partitioned among the
molecules constituting the system.
˙˙˙˙˙˙˙

The partition function provides the most

convenient way for the linkage of microscopic
properties of individual molecules (such as dis-
As the molecules of gas are in constant mo- crete energy levels, moment of inertia and dipole
tion the gas is continuously changing of its own moment etc.) with the macroscopic properties
accord from one microstate to another and almost (such as molar heat, entropy and polarization).
as frequently from one macrostate to another. From Maxwell-Boltzmann’s law, one can
The microstates which are allowed under write
given restrictions are called accessible microstates,
ni g
e.g., in the case of three molecules a, b, c to be = i e–εi/kT
distributed between two halves of a box, if none n0 g0
can be outside the box, then (ab, c), (a, bc), (ac, b) where ni is the number of molecules in ith state, n0
are accessible microstates while (a, b), (a, c), (b, c) is the number of molecules in zeroth level and gi
are inaccessible states. and g0 represents the degeneracies of the ith and
A collection of large number of identical, zero level respectively. When ε0 = 0 then g 0 = 0,
independent system is called an ensemble. In an therefore, one can write
ensemble the systems play the same role as
molecules in a gas. ni
= gi e–εi /kT
n0
Partition Function
or ni = n0 gi e–εi /kT
Statistical thermodynamic analysis is facili-
tated through the use of the partition function Taking summation of this equation, we get
which works as an analytical tool. The partition Σni = n0 Σgi e–εi /kT
function is defined as—
Q = Σgi e–εi /kT Σni = n0 Q
where gi is the statistical weight factor and is n = n0 Q [˙.˙ Σni = n]
equal to the degree of degeneracy, i.e., the number
n
of superimposed energy levels. k is Boltzmann’s or Q =
constant and equal to the ratio of gas constant R to n0
the Avogadro’s number. ε i is the energy of the The partition function may be defined as the
quantum state in excess of the lowest possible ratio of number of particles in the ith level to that
value and T is the temperature on Kelvin scale. of zero level. At absolute zero
In the equation quoted above the summation n → n0 and Q → 1 as T → 0
is taken over all integral of i from zero infinity
corresponding to all possible energy states of the Hence, the value of partition function increases
molecule. with temperature.
384 | Physical Sciences

The Distribution Functions Bose-Einstein Condensation

Consider a system having n particles. These The ideal Bose-Einstein gas is particularly in-
particles are divided into quantum groups or levels teresting because it can undergo a phase transition.
such that there are n1 , n 2 , n 3 … n i … number of When the temperature T of a Bose-Einstein
particles in groups whose approximate constant gas is lowered below the critical temperature T0,
energies are ε 1 , ε 2 , ε 3 … εi … respectively. gi the number of particles in the ground state rapidly
denotes the degeneracy or statistical weight (i.e., increases. The rapid increase in the population
the number of sub-levels) of ith level. Then the of the ground state below the critical
expression for the most probable distributions in temperature T0 for a Bose-Einstein gas is called
the three statistics can be given as— the Bose-Einstein condensation. The critical
(a) Maxwell Boltzmann statistics temperature T0 at which the Bose-Einstein
ni 1 condensation starts depends upon the particle
= α + β εi density n/V of the gas.
gi e
gi
or = eα + β εi
ni ( nn ( 1.0
0
(b) Bose-Einstein statistics
ni 1 0.5
= α + β εi
gi e –1
1.0
gi
or + 1 = eα + β εi O 0.5 T
ni ( (
T0
(c) Fermi-Dirac statistics Fig. represents the fraction of particles
ni 1
= α + β εi condensed in the ground state for T ≤ T0 . At the
gi e +1 ground state E = 0, the particles of a Bose-
gi Einstein degenerate gas condensed in the
or – 1 = eα + β εi ground state do not contribute to the energy.
ni
It is to be noted that if g i/ni is very large in For temperature about T0 , α ≠ 0, the number
comparison to unity, we may write of particles given by
2 ∞ x1/2 dx
gi
~
–
gi g
+ 1 –~ i
V
n = 3 (2πmkT)3/2 ∫ α+x
ni ni ni h π 0 e –1
√
gi with the negligible number of particles in the
i.e., for large value of quantum statistical ground state and the gas is said to be the classical
ni
distribution laws (Bose-Einstein and Fermi-Dirac) or non-degenerate. In the case of Helium in
approach to classical (Maxwell Boltzmann) distri- liquid state T0 can be calculated to have the value
bution law. This is the case of normal existence of 3·12 k. Therefore, the degeneration and
gases when the temperature is not very low and condensation of Helium must start at 3·12 k. But
pressure is not very high. experimental observation, of the lambda point
All the three distribution functions have been transition in liquid helium show that the
compared is the following figure : condensation of helium starts at 2·19k which is
close to the calculated value. Below this
temperature the liquid helium exhibits the
remarkable physical properties of a super fluid. It
is generally believed that the lambda point
transition observed in liquid helium at 2·19 k is
essentially a Bose-Einstein condensation.
Specific Heat of Electron Gas
If the electron gas at absolute zero is heated
the electrons are gradually raised upto higher
states but the change in electron distribution will
 Physical Sciences | 385

at first only take effect at the place where the It is seen at low temperature, the electron
Fermi function falls away and that by slowly specific heat increases linearly as the temperature
rounding off the corners of the distribution curve 3
increases. It approaches the value R at high tem-
as indicated below : 2
peratures.
Planck’s Radiation Law
A black body radiation chamber filled with
ni T=0 simple harmonic oscillators of molecular dimen-
—
gi sions, in addition to radiation. These oscillators
are known as resonators. The energy of resonators
is quantized.
E = nh ν
The average energy of a Planck’s resonator is
εf given by
– hν
The main body at the electrons, however, is ε = hν/kT
left untouched by the rise of temperature. For not e –1
too high temperatures, therefore, only a vanish- The number of resonators per unit volume in
ingly small fraction of the electrons take part in the frequency range ν and ν + dν is given by
the thermal motion, so that the specific heat of 8πν2
electrons is very small. (It is only when high N = dν
c3
temperatures are reached, –~ 104 °C), i.e., for above
The energy density belonging to range dν can
room temperature, that the tight packing of elec- be obtained by multiplying the average energy of
trons is the deeper states gradually becomes a Planck’s resonator by the number of resonators
loosened and we obtain a noticeable contribution per unit volume, in the frequency range ν and
from the electrons to the specific heat of metals. ν + dν, i.e.,
It can be shown that at low temperatures, the 8πν2 . hν
total energy of the electron gas is given by Eνdν = dν
c3 ehν/kT – 1
2 2
E =
3
Nε
5 i [1 + 5π12 (kTε ) + …]
f
=
8πhν3 .
c3
1
ehν/kT – 1
dν
and specific heat Eν dν is the energy density, i.e., total energy
π2 .kT per unit volume belonging to range dν. This law is
C s = Nk. (approx.) known as Planck’s radiation law.
2 εf
The energy density ελdλ in terms of wave-
εf
= γ ′T for T << Tf = length can be obtained by
k
c c
The variation of Cs with temperature is shown ν = , dν = – 2 dλ
λ λ
below :
8πhc . 1
i.e., E λdλ = dλ
λ5 ech/λkT – 1

Specific Heat of Solids

Dulong and Petit’s law—The product of
atomic weight and specific heat (i.e. atomic heat)
of an element at constant volume is constant. For a
number of elements it is equal to 5·96.
But experiments confirms that the law
completely fails at low temperatures.
386 | Physical Sciences

Einstein’s theory—It assumes that atoms of molecules N, so fluctuation is proportional to

solid vibrate simple harmonically with frequency -1
N 2.
(ν) which is same for all atoms and is
characteristic of a solid. For a Grand Canonical Ensemble
dE For an ideal classical (Maxwell's–Boltzmenn)
cν = gas
dt
– (2πmτ)3/2
3Nhν N = eµ/τ V
where E = hν/kT h3
e –1
eθ/T θ 2 – –
∂N
⇒ cν = 3R [ θ/T
(e – 1) 2]( )T ⇒
∂µ
=
N
τ
where N is Avogadro’s number, shows the  –
number of atoms in one gram mole of solid. Each (δN) 2 = N
atom has 3 degrees of freedom. For classical Maxwell Boltzmenn gas
∆N  (δN) 2 1/2
θ =
hν =  – 
where
k N  2 
(N)
θ is called Einstein’s temperature.
 N – 1/2
 1 1/2
Drawbacks—
(i) θ obtained empirically and cannot be
=   ()
 – 2 = N
(N)
verified by any physical data.
(ii) At low temperatures cν ∝ T3 , which this – PV
As N =
theory fails to explain. kT
∆N kT 1/2
Debye’s theory—It assumes that atoms
vibrate in many different modes, the frequency N
= ( ) PV
varying from one to another and the total number For smaller volume of the gas, the greater is
of modes are limited in number. the fractional fluctuation.
Observations— Transport Phenomena
(i) At high temperature cν = 3R is a constant.
● Relative number of gas molecules traversing
(ii) Atomic heat decreases with decrease in
temperature. the distance s without collisions :
(iii) At very low temperature cν ∝ T3. N
= e–s/λ
N0
Fluctuations and Irreversible Processes where, λ is the mean free path.
The root mean square deviation from the ● Mean free path of a gas molecule :
mean is given by 1
–2> λ =
∆n = < (n – n)  2 πd2 n
√
For a physical quantity f, standard deviation where, d is the effective diameter of a
– – molecule, and n is the number of molecule
∆f = [f 2 – (f 2)]1/2
per unit volume.
The probability of distribution decreases and
for large values of n, (n/2, n/2) state is most likely ● Coefficient of diffusion D, viscosity η, and
to occur. heat conductivity χ of gases :
1
Fluctuation is Ensembles D = <ν>λ
3
The energy fluctuation is measured by the 1
ratio (in canonical ensemble) η = < ν > λρ
3
∆E T
= (kCν )1/2 1
E E χ = < ν > λρC v
3
For large values of T, the extensive quantities where, ρ is the gas density and Cv is its
–
C ν and E are proportional to the number of specific heat capacity at constant volume.
 Physical Sciences | 387

● Frictional force acting on a unit area of plates particle moving with a velocity v is –6πRην,
during their motion parallel to each other is a where, R is the radius of the particle and η the
highly rarefied gas : coefficient of viscosity.
1 Therefore, Mα = 6πRη
F = < v > ρ | u1 – u 2 |
6 6kT M
where, u 1 and u 2 are the velocities of the Hence, <r2 > = t·
M 6πRη
plates. kT
● = t
Density of a thermal flux transferred between πRη
two walls by highly rarefied gas : This equation was originally derived by
1
q = < v > ρCv |T1 – T2 | Einstein. It shows that the squared displacement is
6 linear in time. Einstein treated Brownian motion
where, T1 & T2 are the temperatures of the on the basis of the so-called “Random walk
walls. problem”, and thereby established a relationship
between diffusion and the machanism of
Brownian Motion molecular fluctuation.
Small particles when immersed in liquids are Einstein’s relation for mobility—
found to exhibit a random type of motion. This
was first observed by Robert Brown in 1828 and 1
µ =
the phenomenon is known after him as Brownian Mα
motion. kT
As D =
It is now known that the source of this motion Mα
lies in the incessant and random bombardment of is now equal to
the particles by the molecules of the surrounding
liquid. D = kTµ
For Brownian movement This is Einstein’s relation connecting the
6kT 6kT coefficient of diffusion with mobility.
<r2 > = · αT = t If the particle lies under influence of an
Mα2 Mα
6kTτ external electric field E,
= t kT
M D = µ
If the particles are spherical, we can use e
Stoke’s law viz, the viscous drag force on a which is original form of Einstein’s relation.

Solved Questions
Q. 1. (a) What is the probability that a
molecule in the cubical vessel (fig. shown) lies
in the volume ∆ V formed by two planes at z
H dz
and z + dz.
Z

R
dz
Solution : (a) The volume of the vessel is
V = a3 . The volume of the element ∆V = a2dz.
∆V a2 dz 1
(b) Find the probability that a certain Therefore, P (∆V) = = 3 = dz
V a a
molecule lies in a volume element ∆ V formed
by the planes at z and z + dz, in a conical vessel 1
The probability density in this case is ,
whose height is H and the radius of the base is a
R (fig. 2) whose dimension is the reciprocal of the length.
 Unit-6

ELECTRONICS
Electronics Matching
The physics behind the performance of signal In an electric circuit, energy is transferred
processing devices is known as electronics. from the source to various branches in the circuit.
The rate of this energy transfer is known as power
Active Components delivered. The power across the load resistance R i
The active electronic components are those can be given by
which has an ability to amplify the input signal, V20/RL
e.g., transistor (BJT), VJT, FET etc. P =
  Ro  2
Passive Components  1 +  RL
The passive electronic components cannot R L → 0, P → 0
amplify the signal but some of those may change R L → ∞, P → 0
the phase of an input signal, e.g., resistor, capaci-
tor and inductor etc.
Kirchoff’s Laws of Circuit Analysis
Pout
1. Kirchoff’s Current Law (KCL)—At any
instant of time the algebraic sum of current
entering a junction is always equal to zero.
2. Kirchoff’s Voltage Law (KVL)—At any
instant of time the algebraic sum of the voltage RL = RO
RL
drops in a loop to zero.
Model of a Battery For maximum power transfer (dRdP ) = 0
L
We can approximately represent the behaviour or Ro = RL
of a battery as follows :
The maximum power transfer takes place
when the load resistance is equal to the source
resistance. This condition is called ‘Matching’.
Under matched condition the source will deliver
the maximum power
V 20
P max = (for RL = R)
4R L
Hybrid parameters
A battery is consist of an ideal voltage source
and an internal resistance R0 then the voltage at its i1 i2
output terminal will be Input Output V2
V1 port port
V = Vo – iR0
 Physical Sciences | 447

In a two port network an ideal current and an (b) Varactor or varicaps diode—The varac-
ideal voltage source are connected at input an tor is a acronym for voltage variable capacitor
output respectively. So input current and output diode. The junction capacitance of a diode is not
voltage are independent variables, whereas input constant but it a function of applied voltage and
voltage and output current are dependent variables. the design of junction. This variable capacitance
Then the equations relating the variables will be property of diode is useful in tuning of a radio
V1 = h11 i1 + h12 V2 receiver.
i2 = h21 i1 + h22 V2
for V2 = 0
V1 A varactor equivalent
h11 =
i1 If the p-n junction is abrupt, the capacitance
Short circuit input resistance parameter is varies as the square root of the reverse bias. In a
i2 graded junction, the capacitance varies as
h21 =
i1 1 (for hyper abrupt
Short circuit forward current transfer ratio C∝ for Vr >> V0
V nr junction n = 2)
is i1 = 0
V (c) The tunnel diode—This type of diode is
h12 = 1
V2 controlled by a quantum mechanical phenomena
Open circuit reverse voltage ratio is of tunneling between the two sides of a diode.
i An interesting feature of this type of diode is
h22 = 1 its negative resistance region in V-I curve [See
V2
figure].
Since all the parameters have different dimen-
sions, these are known as hybrid parameters. I
Negative resistance region
Special Diodes
(a) Zener diode—The reverse bias break- Ip
down voltage of a junction can be varied by
choice of junction doping concentration. This
breakdown mechanism is called zener effect in
which the covalent bonds in the potential junction
breaks abruptly. A zener is operated under reverse
bias conditions. In a well fabricated zener the Iv
reverse current became independent of reverse
voltage after zener breakdown. That is why a
zener is called breakdown diode sometimes. Vp Vv V

I-V characteristics of a tunnel diode

Zener diode symbol
For negative resistance region the dynamic
dV
resistance is negative. This region is the direct
dI
result of quantum mechanical tunneling between a
degenerate p-type material and a degenerate n-
type material.

Transistors Operating Modes

Depending upon the common terminal
between the input and output circuits of a
Zener V-I characteristics transistor it may be operated in three modes.
448 | Physical Sciences

(a) CB mode—Common base or grounded Crystal Oscillators

base configuration is shown in the figure in which
an n-p-n transistor is connected in CB mode.
IE
+ +
VEB VCB RL
IB VCC
– –
(b) CE mode—In this mode emitter is the
common terminal between input and output mode.
An amplifier when operated in such a mode
provides a phase change of π.
The quartz crystal shows an interesting
property when an alternating electric potential is
applied on its opposite surface. The surface
expands and contracts and resonance takes place if
the natural frequency of the crystal equals the
applied electrical frequency. The crystal frequency
depends on its dimensions. This property of
quartz crystal is used to construct the quartz
oscillator which has precise frequency stability.
A crystal oscillator and its electrical equiva-
lent is shown in the figure alongside. In electric
equivalent circuit the inductance corresponds to
inertia, resistance corresponds to friction and
capacitor corresponds to reciprocal of damping
constant for a mechanical system.

(c) CC mode—In common collector config-

uration the collector is the common terminal as C′ represents the electrostatic capacitance
shown in figure. between electrodes with the crystal as a dielectric.
C′ > C

Operational Amplifiers
The operational amplifiers are high gain
negative-feedback amplifiers which are used to
perform mathematical operations and highly
useful in digital computation.
A schematic diagram of an operational
Comparative performance of CB, CC and CE amplifier is shown below. The common grounded
amplifiers— connection is forbidden.

Parameter CB CE CC

Current gain α β Highest (1 + β)

Voltage gain High Highest Nearly unity

Power gain Moderate Highest Lowest

The negative and positive terminals are
Input Lowest Moderate Highest inverting and non-inverting terminals. The ampli-
impedence fication and feedback are made very large for
extreme stability.
Output Highest Moderate Lowest
impedence The current in the resistors Ri and Rf will be
Ei E
Phase reversal No Yes No and 0 respectively according to Ohm’s law.
Ri rf
 Physical Sciences | 449

By Kirchoff’s first law (KCL) the sum of 1. For the AND of n inputs named A, B, C,
these two currents at junction ‘o’ will be  – – – –
… n, A·B·C … = A + B + C + …… n
Ei Eo
+ = 0 and its dual relation
Ri Rf
– Rf 2. For the OR of n inputs named A, B, C … n
The output voltage Eo = E  – – – –
Ri i A + B + C + … n = A · B · C …… n
The polarity of output voltage is reversed and These laws help in the construction of all the
magnitude is multiplied by ratio of feedback resis- gates from NAND and NOR gates only. That is
tance to input resistance, i.e., Rf / Ri why NAND and NOR gates are known as
The gain can be given by universal gates.
Eo Rf CMOS Digital Logic Inverter
= – (A >> 1)
Ei Ri
The basic circuit element for digital circuit
The feedback amplification is independent of design is the logic inverter. This element is used to
open loop gain if the open loop gain is too high. build logic gates and more complicated digital
Exclusive or Function circuits. The basic CMOS logic inverter is we will
assume the body and source terminals are connec-
The exclusive-OR or XOR is a logic function ted together, so there’s no body effect to consider.
which provides an output only when its both We’ll also assume that the MOSFETs are matched
inputs are at different values. Its symbol and truth
transistors.
table is given below :
VDP
x QP
f=x+y iDP
y
VI Vo
f = x⊕y iDN
x y f
1 0 1
0 1 1 The operation of this circuit can be summa-
rized as—
0 0 0 When V I is ‘low’, QN is ‘off’, QP is ‘on’
1 1 0 ⇒ V0 is ‘high’
The logic operations are when VI is ‘high’, QN is ‘on’, QP is ‘off’
x⊕y = y⊕x ⇒ V0 is ‘low’.
1 ⊕ x = –x Conceptually, this CMOS circuit is intended
to junction as “complementary switches”.
and 0⊕x = x
The directions of the arrow indicate the
–
x.y complementary nature of the two switches. When
x
y one is closed the other is open’s and vice-versa.
f=x+y VDD

Pull up
x . y– o switch

De-Morgan’s Theorems VI Vo
For the simplification of Boolean expressions Pull down
De-Morgan explained the equivalence of OR and switch
AND functions with the help of two laws.
450 | Physical Sciences

Graphical Analysis of the CMOS Logic The vDS values are different for the two
Inverter transistors. By KVL
Here QP is treated as an active load for QN VDD – vSDP = vDSN …(1)
through the converse would produce identical How we interpret this equation can give us
results. the answer on how to plot the QP ch. curve in the
Graphical solution when V1 is “high” iDN – vDSN plane.

VDD The minus sign in (1) tells up to flip the Q P

VSGP characteristic curve about the vertical axis.
QP
VDP VSG = 0 (must be ≥ | V+P | The V DD term tells us to shift this flipped
for induced channel curve VDD units to the right.
i Vo
Graphical solution when vI is “low”—
+
QN VSGP = VDP VDD
VGSN = VDP − QP VSG = VDD > |Vtp| = 0

VGS > Vt ⇒ 0 triode VDP triode or

saturation
or saturation
i Vo
iD
i QN VGS ~
~0
VGSN = VOW = out = VDD + (must be ≥ VtN
VGSN = 0 for induced channel)
−
e
od
tri

int
iD
n i
n

po
Q

ting Load curve Flip & Shift i

era
Op (VSGP = 0) QP Ch. curve to
obtain Load curve

VOL = O VDP V0VDS Operating point

QP off or Saturated
VGSN = VOL = 0 VON ~ ~ VDD

VDD VDS
no power O VDD
i 0 QN off
consumpted Q in triode
or saturated P
−1
Vo = 0
]k (NL(
1
N N
(VDP − Vtn) Ω ] = VDSN VDD
(in triode region) No power
consumed −1
rDSP = ]k (WL( (V
1
P P ]
DD − |Vtp|) Ω
We’ll consider the two extremes of the input : (p in triode region)

VI = 0 (“low”) Vo = VDD

VI = VDD (“high”) (Equivalent circuit)

With QN considered the driving transistor and
QP the active load, then for a graphical solution,
we’ll be plotting characteristic and load curve in
the iDN – vDSM plane.
For this digital inverter circuit in figure we
see that
iDP = iDN [Fig. (1)]
Summary of the CMOS Digital Logic
Inverter
1. The output voltage levels are ~ 10 mV and
~ V DD – 10 mV, thus the output signal swing
is nearly the maximum possible (Gives rise to
the so called noise margins).
 Physical Sciences | 451

2. The static power dissipation is nearly zero (a Where vGS = VDD – vI

fraction of a µW) in both states. The dynamic v0 ≤ vI + | Vtp | 
power dissipation is not zero as the gates are And vSD = VDD – v0
changing states. For the digital logic inverter circuit, these two
3. How output resistance in either state : low currents must be equal—
resistance to ground in the “low” output iDN = iDP …(5)
state, and low resistance to VD D in the
“high” state. To construct the complete characteristic curve
for this logic inverter, we solve (5) for v 0 using (2)
4. Input resistance of the inverter is very large through (5) as vI varies from 0 to VDD.
(ideally infinite). The inverter can derive a
very large number of similar inverters with In the case of symm MOSFETs in which
little loss in signal level. Vtn = | Vtp |
5. The gate input capacitance is not negligible,
so it will take time to “charge” up. Low out-
put resistance helps to keep the time constant
and kn′ (WL) n
= kp′ (WL) , p

τ = Rout Cin small when driving similar the voltage transfer characteristic curve will be
inverter. symmetrical as—

Characteristic Curve for the CMOS V0 QN in saturation

Logic Inverter QP in triode region
QN off A
VON = VDD Slope = −1
Rather using a graphical solution, it is fairly defines VIL
straight forward to numerically, construct the VDD
+ Vt QN &
2 B QP in saturation
characteristics curve v0 versus VI. To obtain an limits for

V 1
es = −
equation to plot for the inverter characteristic both QP & QN QP in saturation

IH
fin e
Saturated C QN in triode region
de lop

curves we will use the equation for the MOSFET VDD

S

− Vt QP off
current in the triodes and saturation region of 2
V
operation. VOL = 0 Vt VIL VIH (VDD − Vt) Vt VDD t
NMH
The iD – v DS characteristic curve for QN is VOH − VIH
VDD
given by Vth =
2

iDN = kn′ ( )[
W
L n
1
(vI – vtn) v0 – v02
2 ] (The voltage transfer Ch. curve)
…(2) NMH = VOH – VIH
vI = vGS 1
NMH = (3VDD + 2Vt)
v0 = vDS 8
v0 ≤ vI – vtn NML = NMH

iDN =
1
k′
2 n L()
W
n
(vI – vtn)2 …(3) Noise Margin
From this characteristic curve we can obser-
v0 ≥ vI – vtn
ves that there is a range of input values (from O to
While the iD – vDS Ch. curve for QP is given VIL) that produces a high output voltages and a
by range of input values (from VIH to VDD) that pro-
iDP =
1
k′( ) (V
2 p L
W
P
DD – vI – | Vtp |)2 …(4)
duces a low output voltages. This is one advantage
that digital circuits have over analog ones.
452 | Physical Sciences

These ranges of input voltages are called 6. IDSS is the constant drain current when VGS
noise margins. In particular, the noise margin for = 0.
high input NMH is defined as— (True for both JFETs and D-MOSFETs).
NMH = VOH – VIH …(7) 7. The devices are square law device.
And the noise margin for low input NM L is 8. Unlike JFETs and D-MOSFETs, the
defined as E-MOSFET can operate with VGS = OV.
NML = VIL – VOL …(8) IDSS
9. Mid-point bias for a JFET is ID = , obtai-
2
Refer to figure above the voltages VIH and ned by setting VGS = VGS(off)/3·4.
VIL are defi-ned where the slope of the curve is
–1. Solving equation (5) at these two points, we 10. MOSFETs differ from JFETs in that the gate
get of a MOSFET is insulated from the channel
by SiO2 layer, whereas the gate and channel
1 in a JFET are separated by a PN junction.
VIH = (5VDD – 2Vt) …(9)
8
11. A depletion MOSFET (D-MOSFET) can
1 operate with a zero, positive or negative gate
and VIL = (3VDD + 2Vt) …(10)
8 to source voltage.
Using (9) and (10) in (7) and (8) we can 12. An E-MOSFET has no IDSS parameter
determine express for two Nm
 VGS 
ID = IDSS  1 –
VGS(off)
1 JFET/D MOSFET
NMH = (3DD + 2Vt)
8

transfer characteristic
NML = NMH
 VGS 
gm = gmo  1 –
=
1
(3VDD + 2Vt)  VGS (off)
8
These two NM are equal as we assumed the 2IDSS
gms =
two MOSFETs in the inverter circuit to be | VGS(off) |
symmetrical. ID = K(VGS – VGS(th))2

Summary for FETS W W

k = kn or kp
L L
1. FET are unipolar devices.
2. The JFET operates with a reverse biased pn Digital Electronics
junction (gate to source).
Binary number system—Base = 2; There
3. The high input resistance of a JFET is due to are two no’s : 0 and 1.
the reverse biased gate-source junction.
● Single binary number is called a Binary digIT
4. Reverse bias of a JFET produces a depletion or bit.
region within the channel, thus increasing
channel resistance. ● Binary number groups called words.

5. For an n-channel JFET, VGS can vary from ● Words are subdivided into 8-bit groups called
zero negatively to cut-off, VGS(off). For a p- bytes.
channel JFET, VGS can vary from positively ● One-half a byte is sometimes referred to as a
to VGS (off). nibble (a term not often used any more).
 Physical Sciences | 453

O16 = O10 = 0000 2 916 = 9 10 = 1001 2 Binary Hex and Hex-Binary Examples
916 = 1 10 = 0001 2 A16 = 1010 = 1010 2 Convert—More binary to hex
B 16 = 11 10 = 1011 2 101110100010; 101101110·01010011;
C 16 = 12 10 = 1100 2 1111111101·10000111
D16 = 1310 = 1101 2 ● To Convert hex-binary
E16 = 1410 = 1110 2 2375]16 =; CD·89; 37AC·6; 3·DCAB.
F 16 = 15 10 = 1111 2 ● Note that leading zeros are added or
removed as appropriate in the conversion
Note—One nibble = One hex digit.
processes.
● Since 24 = 16, each hex digit effectively
represents the same numeric count as four ● Integer Binary/Decimal Conversions
binary digits. Number to be converted :
● Another way to say this is that one column in 1 0 1 0 1 1 0 1
a hex number is the same as four columns of
a binary no. — Power of 2 of column :

162 161 160 16−1 16−2 7 6 5 4 3 2 1 0

1 0 0 1 0 1 0 1 1 0 1 1 0 1 1 1 1 0 1 0 128 64 32 16 8 4 2 1
−6 2−7
210 29 27 6 25 22 20 2−2 2−4 2 128 + 0 + 32 + 0 + 8 + 4 + 0 + 1 = 173
211 2 23
21 2−3 2−5 2−8
28 24 2−1 Fractional Binary/Decimal Conversions—

= 95B·7A ● For numbers right of the binary point, we use

the same approach, remembering that the
Hexadecimal-Binary Conversion powers of 2 are negative

Most computers process 32 or 64 bits at a No. to be converted : 0·1 0 1 1

time.
Power of 2 : –1 –2 –3 –4
Example—
01111000101001011010111110111110 No. represented : ·5 ·25 ·125 0·0625
● Separate into 4-bit groups, starting at the
right. ·5 + 0 + ·125 + 0·0625 = 0·6875
0111 1000 1010 0101 1010 1111
Second example— 0·01011
7 8 A 5 A F
= 0·34375
1011 1110
Decimal → Binary Integer Conversions
B E
Simplify divide the decimal no by 2 :
● Another example—
● The quotient becomes the new no. to divide
Grouping—
again.
1001011100·11110011 2
● The remainder, which is always be 1 or 0,
= 10 0101 1100·1111 0011 becomes one bit of the binary no., least
0010 0101 1100 1111 0011 significant digit first. The last quotient will
= · always be 0, the last remainder always 1.
2 5 C F 3
454 | Physical Sciences

● Div. continues until the quotient is 0 (with a last remainder of 1). Thus, converting 26 to binary.
26 ÷ 2 = 13, remainder 0
13 ÷ 2 = 6, remainder 1 0 110102 = 26 10
6 ÷ 2 = 3, remainder 0 1
3 ÷ 2 = 1, remainder 1 0
1 ÷ 2 = 0, remainder 1 1
1
● Convert 11711 to binary :
Ans. 1110101 2
Converting Decimal Fractions to Binary fractions—
● The method of successive multiplication :
Example 0·25 to binary :
0·25 × 2 = 0·5 : note 0 to left of binary points
0·5 × 2 = 1·0 : note 1 to left of binary points
● There is no part of the decimal fraction left, thus we are done.
● Read the binary fraction as the numbers to left of binary point in the two results, the remainder
from the first multiplication first.
● So 0·25 10 = 0·012
Eq. 0·375 : 0·375 × 2 = 0·75 Multiply decimal fraction by 2. Record the number.
0·75 × 2 = 1·5 left of the binary point. This will always be a 1
0·5 × 2 = 1·0 or 0.
0·37510 – 0·0112 Multiply the remaining fraction by 2, repeating
above action.
0·427 : 0·427 × 2 = 0·854 0·42710 Continue until the fraction is eliminated.
0·854 × 2 = 1·708 01101101 2
0·708 × 2 = 1·416
0·416 × 2 = 0·832
0·832 × 2 = 1·664
0·664 × 2 = 1·328
0·328 × 2 = 0·656
0·656 × 2 = 0·312
Thus we must change the last rule.
Continue until the decimal fraction is eliminated or the binary fraction has at last twice as many
places as the decimal fraction.
Converting mixed decimal numbers—
Read ← for integers
→ for fractions
36·125 = 100100·001 2 ; 19·37510 = 10011·0112
7·33 10 ≈ 111·0101 2
10·17 10 ≈ 1010·00101 2
 Physical Sciences | 455

Signed Arithmetic—In real ‘world’ of n = 0111 1111, then m = 1000 0000 + 1

mathematics, computers must represents both = 1000 0001
positive and negative binary numbers. For integer 2s complement representation
e.g., 124 – 237 = –113 addition and subtraction of both + and – no’s
always work out correctly.
Thus needs to be a consistent methods of re-
presenting negative numbers in binary Arithmetic Furthermore left most bit is always 1 for
operations. negative no. always 0 for a +ve number.
There are various approaches, by they all This means that an n bit binary number in
involve using one of the digits of the binary no to two’s complement can represent a magnitude
represent the sign of the number. range of ± 2 n – 1 – 1.
Two methods are the sign/magnitude repre- Eq. Let n = 0000 0000
senstation and the one’s complement method of —
representation. Then m = n +1

Binary Sign Representation— = 1111 1111 + 1

Sign-magnitude—The left bit is the sign (0 = (1) 0000 0000

for +ve numbers and 1 for – numbers) all are the (1) discard, 0000 0000.
number magnitude. The 2’s complement 0 is 0.
1. _______ Example : –50 : 50 = 0011 0010;
left bit is the sign bit 1’s C = 11001101;
Advantages— 2’s C = 11001110
Simple to implement. –127 : –127 = 0111 1111;
Useful for floating point representation. 1’s C = 1000 0000;
Disadvantages— 2’s C = 1000 0001
Sign bit independent of magnitude; can be –1 1 = 0000 0001;
both + 0 or –0.
1’s C = 1111 1110;
One Complement—
2’s C = 1111 1111
The negative numbers of the same magnitude
But +ve decimal no’s are converted simply to
as any given positive numbers is its one’s comple-
positive binary no’s as before (no 2’s comple-
ment.
ment).
—
If m = 01001100, then m = 10110011
Example : + 67 → 0100 0011
The most significant bit is the sign is 0 for
2’ complement Binary to decimal (2)
+ve binary no and – for negative number.
— 1111 1111 → 0000 0000 + 1
Note the problem : If m = 00000000, then m =
1111 1111; there are two zeros in the method. = 0000 0001
2’s Complement Negative Binary Nos. = 1 → –1
—
m is the 2’s complement of n, then : m = n 1010 0011 → 0101 1100 + 1
+1
= 0101 1101
Eq. n = 01010100, then m = 10101011 + 10 =
1010 1100 = 93 → –93
But for +ve binary no
n = 010111 11, then m = 10100000 + 1 =
1010 0001 0000 0001 → rest a negative no → 1.
456 | Physical Sciences

2’s complement Binary Math Excess –3 code

If we represent binary no’s in 2’s complement To add in excess –3 add the binary no
form, simple addition and subtraction is easy.
8 7
To a – b = a + (2’s complement of b)
— +3 +3
= a+ b +1
Adding no’s is the same as decimal addition 11 10
Subtract 0111 0101 from 0111 1100 : ↓ ↓
2’s complement of 0111 0101 is 1000 1010 1011 1010
+ 1 = 1000 1011. ASC II code (2)
Add : 0111 1100 124 There are 128 basic ASC II characters 0
1000 1011 –117 – 127 10, or 0 – 7616 or (0000 0000 to 0111 1111
binary).
Thrown away ← (1) 0000 0111 007 Gray-Code—Called min. change code t and
Add : 1100 0001 + 0110 1110 is unweighted code. Because of this gray code is
not suitable for arithmetic operation but finds
Note that 2’s complement of 1100 0001 is application in input-/output devices, some types of
0011 1111 so first no is –63 analog-to-digital converters; designation of rows
So 1100 0001 –63 and column in karnaugh map.
0110 1110 110 A three 3-bit Gray code may be obtained by
merely reflecting the two bit code about an axis at
Thrown away ← (1) 0010 1111 + 47 the end of the code and assigning a third bit as O
above axis and 1 below the axis,
Subtract : 1101 1101 from 0101 1100 00 000 0000
(We are subtracting a negative no) : 01 001 0001
0101 1100 92 11 011 0011
+ 0010 0011 –(–35)
10 010 0010
2’ complement 0111 1111 + 127
(a) 110 0110
Add : 1000 0001 + 0111 0010 : 111 0111
1111 0011 → negative number 101 0101
114 100 0100
–127 1100

–13 1

Check 2’s C of 1111 0011 → 0000 1100 + 1 1

= 0000 1101 1

= 13 ∴ Binary code.
So –13. 110 111
11
10

11

0
1

Binary codes— 101

001 010 100
BCD code (8421) : 5 6 210
00 01
0

00
0
0

11
11

010
0
00

0101 0110 0010

1
 Physical Sciences | 457

From going 7(0111) to 8(1000) require four Finally more complicated

digits change but Gray code requires (0100 to —
(a + b)·( c + d) = f
1100) only one bit change convert binary no’s
101102 to Gray code a
b
1 0 1 1 02 1 0 1 1 02 Binary f
c
Gray d
0 1 1 1 1 0 1
(Discard Gray)
Any Boolean function f maps a total of n
inputs (x1 , x2, x3, … xn ) into one output (0, 1), that
Convert Gray to Binary
is a Boolean function f of n variables produces a
Convert Gray to Binary single output for each unique combination of
1 1 0 1 1 inputs.
Principle
1 0 0 1 02
discard carry Any Combinational Logic Function may be
represented by two levels of logic—One level of
Logic Gates and Boolean Algebra AND gates and another level of OR gates.
SOP and POS Boolean Representation—If
AND, OR, NOT, NAND, NOR, XOR, we have AND function followed by an
XNOR gate etc. OR, then we have a “Sum Of Products” SOP
Creating Boolean Functions representation.
By Combination of Boolean expression. If we have OR’s followed by an AND,
Boolean functions created in this way are Product Of Sums (POS) representation.
called combinational logic and they are time- A sum of products function looks like
independent. f = (a·b) + (c·d)
Let’s now examine Boolean function → f = (a·b) + (c·d)
circuit design process.
a
Consider these function
b
a f
c
b f d OR level
c
AND level
a+b+c=f
a Similarly POS :
b f f = (a + b)·(c + d)
c
— Using a Truth table to find Boolean
a+ b =f expression—Making a truth table to represent the
a input/ output.
b f
Defining a Boolean expression in SOP or
(a·b) + c = f POS form that satisfies the truth table.
a Constructing a circuit that represents the
b “stuff in the middle” that performs the Boolean
f
function on the inputs to get the desired output

Defines Boolean Design digital

Establish Construct exp in Circuit based on
in/out T.T. SOP or POS Boolean exp
458 | Physical Sciences

Example SOP 0 1 0 0 x + –y + z
Specification : 0 1 1 1
x y z f AND’s 1 0 0 1
0 0 0 0 1 0 1 1
0 0 1 0 1 1 0 1 –x + –y + z
● Output is true (= 1) 0 1 0 0 1 1 1 1
–
For two Combination 0 1 1 1 — x yz In the SOP example, we looked for the output
– in F.
Of the inputs 1 0 1 1 x yz
If we make the 1 1 0 0 For POS, we look for outputs that are 0.
Boolean expression Here, f is 0 only for two combinations of x, y
and z.
SOP, 1 1 1 0
● We represent these conditions in OR form.
The first level will
● Each OR will produce a 0 only for the exact
have only two
– – OR expression shown.
So, f = ( x·y·z) + (x y·z)
f = (x + y– + z)·(x– + –y + z)
AND terms
● The Boolean expression is now. For
● We then define two
f = (x + y– + z)·(x– + –y + z)
AND functions that result in 1’s being output
Note that f = 0 fully defined.
function 1 for each x, y, z combination.
SOP Versus POS Representations
Boolean Notation; Minterms and Maxterms
● Any Boolean expression can be represented
● SOP is the more common form of notation.
in SOP of POS form.
Since SOP expressions can contain a number
of AND terms ORed together, the AND a b f SOP POS
symbol and the parentheses are often omitted 0 0 0 a+b
for simplicity. –
0 1 0 a+b
● Thus, the SOP expression –
f = (a·b·c) + (a·c·d) + (a·b·d) + (b·c·d) 1 0 1 ab
could also be written as : 1 1 1 ab
f = abc + acd + abd + bcd ● Consider the truth table shown.
● The AND symbols are implied in this ● If we consider f in SOP form,
–
representation and are still used when the We have f = ab + ab
expression might not be clear if they were left ● And in POS form
out. –
f = (a + b)·(a + b)
● Since SOP is more compact, it is some times ● Either representation is correct. In general,
called the minterm form, and the AND com- we pick the representation that is simplest.
ponents minterms. Similarly, POS is often ● For example
referred to as the maxterm form, and the OR’s
maxterms, since POS terms are more x y z f OR′s AND′s
unwildly. 0 0 0 0 x+y+z
● The names minterm and maxterm will be 0 0 1 0 x + y + z–
occasionally in future. 0 1 0 0 x + –y + z
● Product of Sums Boolean Form 0 1 1 1 ………… –x·y·z
A POS representation will contain a first level 1 0 0 0 x+y+z
of OR’s.
x y z f 1 0 1 1 ………… x·–y·z
0 0 0 1 1 1 0 0 –x + –y + z
0 0 1 1 1 1 1 0 –x + –y + –z
 Physical Sciences | 459

POS : Boolean Theorems

f = (x + y + z)·(x + y + –z) ● The Boolean functions AND and OR are
commutatives.
·(x + –y + z)·(x– + y + z) x·y = y·x
·(–x + –y + z)·(x– + –y + –z) And x+y = y+x
SOP : ● The Distributivity, the AND and OR func-
f = –xyz + xyz
– tions are both distributive :
● Finding the truth table, given SOP x·(y + z) = (x·y) + (x·z)
Expression and x + (y·z) = (x + y) · (x + z)
We can illustrate the AND version of this
Boolean expression-truth table. relation as—
● Consider the SOP Boolean exp.— x
– + xyz– + xyz– x.(y + z)
f = xyz + xyz y
● We insert the 1 outputs in the truth table and z
By default all the other values of f must be 0. x
x y z f y
= (x.y) + (x.z)
0 0 0 0 x
0 0 1 0 z
0 1 0 0 Similarly OR version can be illustrated.
0 1 1 1 –xyz Other Fundamental Boolean Theorems
1 0 0 0 OR Form AND Form
1 0 1 –
1 xyz x+0=x x·1 = x
1 1 0 1 xy–z x + –x = 1 x·–x = 0
1 1 1 1 xyz x+x=x x·x = x
Step-by-Step Logic Circuit Design x+1=1 x·0 = 0
# And, Finally, we note that :
● We usually go truth table-circuit however, –
consider truth table circuit x=x
x y f SOP term De Morgan’s Laws—OR and AND versions
of the law.
0 0 0
● The OR version of De Morgan’s law states
0 1 1 –xy that the complement of (x OR y) is equal to
1 0 1 xy– x-complement and y-complement, or—
1 1 0 This may also be
stated as—
x  – – – –y
x + y = x·
x + y = x· y
y ● Likewise, the AND version states that the
xy + xy complement of (x and y ) is equal to x-
x complement OR y-complement, or—
y 
x ·y = –x + –y

Looking for the Simplest Expression (also stated as x·y = –x + –y)
● SOP and POS Boolean expression may NOT Simplifying a Boolean Expression—
be in the simplest form : Consider f = abc + abc–
● However, it is possible to reduce a Boolean SOP form, [distributive]
expression to its simplest form using the f = ab(c + c)–
simple identities (we will study simplification = ab(1) = ab
of a circuit using a graphical technique called (Equivalent expression much simpler, there-
the karnaugh map). fore easier and cheaper to build)
460 | Physical Sciences

More complicated example Assume Morgan’s law or other Boolean identities to

f = (((abc + ab + bc)·a + b)·c) + abc simplify an expression, if simplification is
NOT in SOP or POS form. possible.
Using the truth table, the equivalent SOP Exclusive OR/Exclusive NOR (XOR/XNOR)
expression in— ● There are useful logic functions.
f = –abc + abc ● Both have two or more inputs.
= (a– + a)bc ● (Two input truth table shown).
= bc a b a XOR b a XNOR b
0 0 0 1
b′ because –a + a = 1 0 1 1 0
a b c f SOP 1 0 1 0
0 0 0 0 1 1 0 1
0 0 1 0 ● For > 2 inputs, the XOR output is 1 for an
0 1 0 0 odd number of 1 inputs; XNOR has a 1
0 1 1 1 –abc output for an even number of 1 inputs.
1 0 0 0 XOR and XNOR are not basic Boolean func-
1 0 1 0 tions, but can be made from AND, OR and
1 1 0 0 NOT
1 1 1 1 abc a
XOR a b
b
Implemented—
XNOR a
a a b
b
XOR = ab + ab
a
b
(ab + ab)
f
c
b
Decoders
● An n – to –2n decoder is a combinational
a logic circuit that has n inputs and upto 2n
output. That is, it can have 2 n outputs, but it
b
f f may have less.
c
Simplified Simplified SOP
● Each output of a decoder will normally be
c form true (i.e., go to logic 1) for only one
Proper SOP form combination of the n inputs.
● Consider the case of an n = 2 decoder. The
Summary
decoder will have 2 inputs and upto 2n = 22
● Digital circuits are electronic implementations
= 4 outputs.
of Boolean expressions.
● Assume that the decoder has the maximum
● Computer circuits are all composed of logic possible number of outputs (4). Then the truth
gates (AND, OR, NOT, NAND, NOR etc.) table for the 2-input decoder will show that
● Logic circuits whose outputs only depend on for each combination of y and x (00, 01, 10,
their inputs are called combinational logic. 11), one of the outputs will go high (logic 1).
● Combinational logic with more than two ● Let us call the inputs y and x and the outputs
levels of logic may always be simplified into a, b, c and d (here, y is the more significant
SOP or POS form using a truth table. bit). Then let us define a = 1 for y = 0, x = 0;
● An SOP and POS form may not always be in b = 1 for y = 0 and x = 1 c = 1 for y = 1 and
the simplest possible form. We can use De- x = 0 and d = 1 for y = 1, and x = 1.
 Physical Sciences | 461

Truth tables for a-d in our 2 to 4 decoders The multiplexer is usually symbolized by the
are— abbreviation MUX as the symbol for its function.
Input Output Multiplexer : An “Input to Output” selector.
y x a b c d The truth table for a multiplexer is shown
0 0 1 0 0 0 below.
0 1 0 1 0 0 Assume a 4-input MUX, with inputs labelled
1 0 0 0 1 0 a, b, c, d.
1 1 0 0 0 1 Then there must be two address lines, y
⇒ Using the truth table above, we can define (MSB) and x(LSB).
a – d in terms of x and y The output is denoted as f
a=yx
y x f
b=yx Boolean exp. for 0 0 a
= a, b, c & d in 0 1 b
c=yx terms of x & y
d=yx 1 0 c
1 1 d
Now circuit a
x
a
b

b f
y c Output

c d

d x Selector

In the same way, a decoder with 3 inputs will

have up to 23 = 8 outputs and each outputs y
will have a unique address that represents one
of the eight possible combinations of the three
inputs.
Like wise, a 4-input a
3-input b
decoder would have upto x eight c Decoder
24 = 16 outputs, each of output d Note—That the multiplexer has 1-bit output.
which has a unique y decoder e
address in terms of the
z
f Differences in Decoder and Multiplexer
four inputs. g
Decoder—A decoder has n inputs, which are
Note again that any n-p h
called the address.
n
input decoder can have upto 2 outputs, A decoder has up to 2n outputs (it can have
but it may have less. that many, max, but it might have less). Each
output line is true (or 1) for a specific combination
Multiplexer of the input lines, called the address.
Definition of a Multiplexer—A multiplexer Multiplexer—A multiplexer has two sets of
is a combinational logic circuit that has upto 2n inputs—n address lines (fuse like decoder) as
inputs; an n-bit address and one output. many as 2 n inputs, one of which is selected by
The multiplexer connects one of the inputs to each address for outputs.
the outputs, depending on the value of the n-bit A multiplexer has only one output. The
address. output is the value of the input selected by the
The n-bit address is decoded, just as before. address.
Thus the multiplexer uses a decoder and Thus we see that a decoder makes up a part
selector circuit to tie one of its inputs to its output. of a multiplexer.
462 | Physical Sciences

Binary Arithmetic Circuits—A binary adder Thus, to add two 8-bit numbers, we simply
is a large part of a computer CPU. The CPU use eight column adder.
“figuring” unit is sometimes called the ALU, or Alternate adder—
arithmetic/logic unit. a
The ALU is a combinational logic unit that S
can add subtract or do logical operation such as
AND, OR, NOT etc. C0
b
Consider the “rules of addition” (adding two
no’s only). The “Full Adder”—The half adder circuit is
No’s are added on a Columnar basis, starting an important step in building a digital adder, but it
or the right. is not the total solution to the binary addition
If the sum of one column is a 2-digit number, problem.
the right no goes in current column and the left no. When adding a binary no’s on a columner
(always a 1) become a “carry” to the next column basis, in addition to the column bits ai and bi, there
to the left. is also the possibility of a carry bit ci – 1 from the
Column addition always includes the carry column to the right.
from the column to the right. Col (i + 1) Col (i) Col (i − 1)
The resulting sum in any column will be only ci ci −1 ci −2
a single digit. + + +
ai +1 ai ai −1
Binary no. added in same way— +
Basic principles for an n-column addition. bi +1 bi bi −1
Add column i plus carry from column i – 1 si +1 si si −1
(ci – 1)
Thus the 2 bit column adder must be a 3-bit
1 1 1 1 1 1 11
Carries 1 1 1 Carries adder.
1 0 1 010 0 0 1 1 1
5 8 7 This 3-bit adder is a full adder.
6 4 3 0 1 1 111 1 0 1 0 1
In adding two n bit binary number, the right-
(1) 2 3 0 (1) 0 0 1 0 01 1 1 1 0 0
most bits of the two no’s of two number could be
added with a half adder (no carry in).
0
If a one digit result, that number is the sum of However most adder circuits use all full-
column i, si. adder, since carry in for the least significant bit
can be useful we will see that use later.
If a 2-two digit result, the right digit is si : the
left digit is i or ci. Full Adder
c0
The column i carry, c i, will be added to Exp.
column (i + 1). a b ci s c0 s
A Half Adder 0 0 0 0 0
1 0 0 1 0 abc
a b s c0 eexp. c0 exp.
0 1 0 1 0 abc
0 0 0 0 1 1 0 0 1 abc
0 1 1 0 –ab 0 0 1 1 0 abc
1 0 1 0 – 1 0 1 0 1 abc
ab 0 1 1 0 1 abc
1 1 0 1 — ab 1 1 1 1 1 abc
a –– – – –– ––
S = ab c + abc + a bc + a bc + abc
– – –
Sum
C 0 = abc + abc + abc + abc
b Alternate constructions—
For the full adder
S = (a ⊕ b) ⊕ c
Carry out
C = a·b + (a ⊕ b)·c
 Physical Sciences | 463

Where C = C in Similarly 4-variable—

Cin yz yz yz yz
S
0000 0001 0011 0010
wx
0 0 1 1 3 3 2 2
0100 0101 0111 0110
wx
a Sab 4 4 5 5 7 7 6 6
SOP K map.
1100 1101 1111 1110
wx
ab C0 C 12 D 13 F 15 E 14
b 1000 1001 1011 1010
wx
Simplifying Logic Circuits with Karnaugh 8 8 9 9 B 11 A 10
Maps— Karnaugh Map Comments—
– – k-maps can be labeled many ways, by
f = abc + abc + abc
We use graphical technique to solve this We can label it in binary, decimal, or
which is called Karnaugh map. k-map. hexadecimal. We can also designate the Boolean
function by the k-map squares it occupies.
A k-map is a truth table graph, which aids in
visually simplifying logic The minterms on the k-map can be labelled as
f = Σm(5, 7, 13, 15) in decimal, or
It is useful for up to 5 or 6 variables, and a
good tool to help understand the process of logic f = Σm(5, 7, D, F) in hexa decimal
simplification. yz yz yz yz
0000 0001 0011 0010
A 2-variable k-map. wx
This is very simple k-map demonstrates that 0 0 1 1 3 3 2 2
an n-variable k maps contains all combinations of 0100 0101 0111 0110
wx 1 1
the n variables in the k-map space. 4 4 5 5 7 7 6 6
y y 1100 1101 1111 1110
00 01 wx 1 1
x This minterm express C 12 D 13 F 15 E 14
0 1
10 11 by f = xy 1000 1001 1011 1010
x wx
2 3
8 8 9 9 B 11 A 10
3 variables—Each square represents 3- If we were given the sigma (Σ) co-ordinates
variable minterms or maxterm. above, we could immediately deduce that SOP
All the 8-possible 3 variables terms are function was
represented on the k-map. – × yz
– + wxyz
–
f = w + wx–yz + wxyz
y z y z y z y z
Any square or rectangular group of cells that
000 001 011 010
x 0 1 3 2
is a power of 2(1, 2, 4, 8, 16) is called implicant.
100 101 111 110 All of the groups of squares in the k-map to
x 4 5 7 6 the left (including the single square) represent
As an example, this minterm implicants of different sizes.
cell (011) represents the minterm 0000 0001 0011 0010
f = xyz 0 0 1 1 3 3 2 2
When moving horizontally or vertically, only 0100 0101 0111 0110
1 variable changes between adjacent square never 4 4 5 5 7 7 6 6
2, This property of k-map, is unique and accounts 1100 1101 1111 1110
for unusual numbering system.
C 12 D 13 F 15 E 14
The k-map shown is one labeled for a SOP
1000 1001 1011 1010
As an example, this minterm cell (011)
represents the minterm f = –xyz. 8 8 9 9 B 11 A 10
464 | Physical Sciences

Prime Implicants— ⇒ f = xz–

yz yz yz yz y z yz yz y z
wx wx
wx 1 1 wx 1 1
wx 1 1 wx 1 1
wx wx
So, the prime implicants at right may be
represented by the Boolean expression A POS k-Map—
f = xz The procedure is the same, except that we
k-map as follows— map 0’s
Let – + –x + –y + –z)·(w
f = (w – + –x + –y + z)
For an SOP expression, find the k-map cell
for each minterm of function f, and place a one (1) – + x + –y + –z)·(w
·(w – + x + –y + z)
in it. Ignore 0’s. Simplified expression is
Circle groups of cells that contains 1’s – + –y
f = w
The no. of cells enclosed in a circle must be a
w
power of 2 and square or rectangular. ⇒ f
y
It is okay for groups of cells of overlap.
Each circled group of cells corresponds to a y+z y+z y +z y +z
prime implicant of f. 0000 0001 0011 0010
w+x
yz yz yz yz 0 0 1 1 3 3 2 2
0100 0101 0111 0110
wx w+x
4 4 5 5 7 7 6 6
wx 1 1 1100 1101 1111 1110
w+x 0 0
wx 1 1 C 12 D 13 F 15 E 14
1000 1001 1011 1010
wx w+x 0 0
8 8 9 9 B 11 A 10
Note that the more cells a given circle
enclose, the fewer variables needed to specify the Summary of Karnaugh Map Procedure—
implicant, (1) Start with the truth table of Boolean
– + xz
f = wx expression.
(2) If starting from the truth table, write the
y z yz yz yz
Boolean expression for each truth table term that
wx is 1 (if SOP) or 0(if POS).
wx 1 1 1 1 (3) Develop the full Boolean expression, if
necessary, by ORing the AND terms together if
wx 1 1 SOP or ANDing OR terms is POS).
wx (4) On the k-map, plot 1’s (for SOP) or 0’s
(for POS).
Another example— (5) Group implicants together to get the
The implicants of f are 4, 6, 12, 14. largest set of prime implicants possible. Prime
implicants may overlap each other. They will
This corresponds to
always be square or rectangular groups of cells
– y––z + wxy
f = wx – –z + wxy––z + wxyz– = 0 that are power of 2(1, 2, 4, 8).
 Physical Sciences | 465

(6) The variables that make up the term (s) of y z yz yz yz

the new expression will be those which do not
wx 1
vary in value over the extent of each prime
implicant. wx 1 1
(7) Write the Boolean expression for each wx 1 1
prime implicant and then OR (for SOP) or AND
(for POS) terms together to get the new wx 1
expression.
Example—Reduce the Boolean expression f = xyz + –yz
– yz
f = wx– + wxyz
– – –z + wxyz
+ wxy – 1

+ wxyz + wxy– –z + wx– y–z + wxyz

–
1 1
Ans. f = xy + xz + wz 1
1
Original Circuit—
1
w
–
f = xz + yz
x Prime Implicant should overlap, if this mean
E that they can be made larger resulting gate
Simplify
x
y
z f
z
y
x The simplified expression (and not very sim-
f plified at that) is—
y – +w
f = wz – –xy + x –yz + –xyz
(3) Biggest Prime Implicant Sometimes Major Simplification is Not
Possible—Using the Σ notation
w x y z f
f = Σm(1, 2, 3, 5, 7, 11, 13)
0 0 0 0
0 0 0 1 1 y z yz yz yz
0 0 1 0 wx 1 1 1

0 0 1 1 wx 1 1
0 1 0 0
wx 1
0 1 0 1 1
0 1 1 0 wx 1

0 1 1 1 1 The concept of “Don’t Cares”—The six

1 0 0 0 implicates on the k-map shown be represented by
1 0 0 1 1 the simplified exp. f = xy + xz
y z yz yz yz
1 0 1 0
1 0 1 1 wx 1 1 1

1 1 0 0 wx 1 1 1
1 1 0 1 1 wx 1 1 1
1 1 1 0
wx
1 1 1 1 1
466 | Physical Sciences
y z yz yz yz
x necessarily change when an input changes, but is
wx
wx x 1 1 1
wx x 1 1 1
wx x Representation of sequential logic
Suppose for our particular logic system, we
know that y and z will never be 0 together.
Then the –y · –z implicants do not matter, since
they cannot happen.
The show the condition cannot occur, we put
X’s in the column-they are “don’t cares”—
conditions that cannot happen.
Since the function will never get into the left
column, we “don’t care” how those minterms are
represented.
Why not make them 1’s—Doing that, we
can make a much larger prime implicant and a
simpler Boolean expression. For this implicant,
f = x.
The expression is valid, since the forbidden
condition with never allow the two left squares to
be occupied.
“Don’t Cares” let us further simplify an
expression.
“Don’t Cares”—Summary—“Don’t Cares”
occurs when there are variable combinations,
represented by squares on a Karnaugh map, that
cannot occur in a digital circuit or Boolean
expression.
Such a square is called a “Don’t Care”. Since
it can never happen, we can assign a value of 1 to
the square and use that 1 to (possibly) construct
larger prime implicants, further simplifying the
circuits.
We will see later in sequential logic that the
concept of “Don’t Cares” will help in simplifying
counter circuits.
Sequential Logic and Clocked Circuits
A combinational logic circuit operates such
that the output of the circuit depends only on the
input (s).
Sequential logic differs from combinational
logic in several ways—
⇒ Its outputs depend not only on logic inputs
but also the internal state of the logic.
⇒ The output of sequential logic does not
synchronized to some “triggering event”.
⇒ Sequential logic is often synchronized or
triggered by a train of regular pulses on a serial
input line, which is referred to as a “clock”.
Clock or State or Combinational
Timing Memory Logic
device Element Elements
Output
Input variables
Sequential logic has inputs and outputs like
combinational logic. It also has a timing (synchro-
nizing) mech. and state or memory elements).
Sequential logic may (usually will) have
combinational parts.
Sequential logic outputs depend on the timing
and state elements as well as the input variables.
That is, the outputs normally change as a function
of the timing element.
The Simple Latch or R-S Flip-Flop—The
simplest example of a sequential logic device is
the R-S Latch or R-S flip-flop (R-S FF).
Set
Q
Q
Reset
Rs Flip-Flop
This is a non-clocked device that simply
consists of two cross-connected 2-input NAND
gates (it may also be represented as combina-
tions of other gates).
Note that per our definition, the R-S FF has
an output that depends on its current state as well
as the inputs.
This most basic form of the R-S FF uses
“negative-true” input logic. This is, it changes sta-
tes based on inputs that transition to 0 not to 1.
—
When the Q output is 1 and Q = 0, the RS-FF
is said to be “set”.
—
Likewise, when Q = 0 and Q = 1, the flip-flop
is said to be “reset”.
—
Q and Q will always have opposite states.
The R-S FF is bistable. That is, it is “at rest”
in one of its two states (“set” or “reset”) until an
input forces it to change.
 Physical Sciences | 467

If the RS FF is in the “set” state, it will not go Note the race condition that is triggered by
“reset” until the Reset line goes “true” (in this R = S = O—
—
case, to 0). Then Q = Q = 1, so both are 2-NAND inputs
Likewise, when “reset”, it will not go “set” are 1.
unless the set line goes to 0.
If S and R go to/simultaneously, then all 4
Note also that once “set”, if set goes to 0 inputs of the two 2-input NAND gates are 1 and
more than once, the FF simply stays “set”. both outputs go to 0.
Likewise, when “reset”, more Reset’s do not
The result is a “race” to see which output gets
affect the circuit; it remains “reset”. to 0 first, getting one 2-NAND input to 0, and
Thus the RS FF has an output that depends therefore forcing the NAND output to 1. The
not only on the inputs but the current state. result of the “race” cannot be predicted.
Note—The triggering signal of an input (set Thus R and S = 0 together is forbidden, since
or reset) is always assumed to be momentary. the output is not stable.
That is, a “Set” or “Reset” signal is considered to
as a very short time (in case of real circuits ~ RSFF—Many exotic flip-flops are based on
nanosecond or less). the simple R-S FF.
Assume the FF is Reset (Q = 0). Also, since Used in
the Set and Reset inputs are not active, both input The data storage register.
are at 1. The binary counter.
—
Thus set = 1, reset = 1, Q = 0, Q = 1. The shift register.
Then the cycle is— Most “status indicators”
1. Set goes active (Set → 0).
Many types of computer memory.
2. Then Q must → 1 (output of a NAND = 1
if any input = 0). Note—Q can be set to 1 or 0. Thus the output
Q can be said to “remember” are bit.
3. Then both input to bottom NAND are 1
— When an RS FF has its Q output changed to 1
and Q → 0.
or 0, the output stays in that state until the
4. The other input to the upper NAND is now opposite input is triggered.
0. Thus, when the set signal goes back high, Q
remains at 1, since the other input is still 0. Thus the flip-flop or later has the property of
“remembering” a one-bit binary no. It can be set
5. Likewise, since both inputs to the lower to 1 or 0.
—
NAND are now 1, then the value of Q remains 0. The flip-flop therefore becomes one storage
6. The reverse cycle (set-to-reset) occurs in memory choice i.e. we can use flip-flops to store
the identical way, except that the change is binary no’s in a computer.
initiated by reset going low.
Simplified Version
The truth table for RS FF
S R Q
Inputs Current Outputs New Outputs
— — 0 0 Indeterminate
S R Q Q Q Q
0 1 Set (1)
1 1 1 or 0 0 or 1 Same Same
0 1 0 1 1 0 1 0 Reset (0)
0 1 1 0 Same Same 1 1 Hold
1 0 1 0 0 1 RS FF using NOR gate
1 0 0 1 Same Same 1
0* 0* 1 or 0 0 or 1 1* 1* Q
This is race condition and not stable. The S-R
input 0 -0 is forbidden.
As noted on the t-t, 0-input to both R and S is Q
forbidden. 2
468 | Physical Sciences

Simplified Version Equation the clock on next page does not

S R Q have a 50/50 duty cycle. It stays in the “1” state
0 0 Hold about 35% of time, and in the “0” state about 65%
of the time. Thus
0 1 Reset (0)
1 0 Set (1)
1 1 Intermediate
T T
Inputs Current Outputs New Outputs
— Thus we say that the clock has a “35/65” duty
S R Q Q cycle. In the same way, a clock can have a 70/30
1 2
0 0 1 or 0 0 or 1 Same Same cycle time (i.e., it stays in the “1” state 70% of the
time) and so on.
0 1 0 1 0 1
0 1 1 0 Same Same Clocked Flip-Flop
1 0 1 0 1 0
1 0 0 1 1 0 All ff’s have the same basic configuration—
1* 1* 1 or 0 1 or 0 0 0 Both true and false outputs (“Q” and “Q-not”)
Intermediate state “Set” is when Q = 1.
Triggered by “Set” and “reset” input.
Clocked The most useful ff’s are not simple
The majority of all sequential logic circuits asynchronous (non-clocked) ff’s, however, but
are clocked logic circuits. synchronous (“clocked”) ff’s.
Clocked circuits are circuits that are regulated
by a clocking element, which determines when Clocked ff’s are very similar to non-clocked
state changes are made in the sequential logic ff’s the main difference is that in addition to a
circuit. “set” or “reset” input to cause the outputs to
change, there must also be the presence of a clock
In a clocked sequential circuit, in general, signal in its true state (normally 1).
the circuit can only change states on a “tick of the
clock element. Thus we say that all operation in Thus clocked ff’s do not change states,
this type of circuits are “clocked”. regardless of the set or reset inputs, until the
Clock does not mean any kind of watch. “clock ticks”.
We refer to a circuit as a “clocked circuit” Clocked RS Flip-Flop
when all the sequential elements in the circuit In addition to the set and reset inputs, the
change states in synchronization to a train of clock input is present.
pulses.
Such a pulse train is shown below. The clock Set +
Q
pulse change regularly from 0 to 1 and back. S
1 Clock
R
0 Q
Time Reset +
T T
(NAND Version)
Facts—The period T of a sequential logic
clock is the distance between two identical points Since when clock is low (0), neither set or
on the pulse train, e.g., two rising edges or two reset affect the circuit, we say that the clock
falling edges. “gates” the set or reset signal to the RSFF.
The clock has only two states : 0 and 1, and In this case, the set or reset input must be high
clock alternate between these states. (1) to set or reset the ff when the clock goes true
The amount of time that the clock spends in (0 → 1).
each of the two states is called the duty cycle. Having set and reset 1 at the same time is
The clock above has a 50/50 duty cycle; it forbidden as for the RSFF, simultaneous set and
spends equal time each period in the 1 and 0 reset true causes a race condition when clock is
states. high.
 Physical Sciences | 469
—
Clock S R Q Q New Output There are many ff’s, (Summary), the RS FF
0 X X 1 or 0 0 or 1 Q Q
— is the simplest example of sequential logic.
Same Same (1) Non clocked; made of crossed-connected
1 0 0 1 or 0 0 or 1 Same Same 2-input NANDs. (or NOR’s etc.)
1 1 0 0 1 1 0 (2) Output depends on current state as well
1 1 0 1 0 Same Same as the inputs.
1 0 1 1 0 0 1 (3) Uses “negative-true” input logic. Output
—
1 0 1 0 1 Same Same Q and Q are always opposite.
—
1 1* 1* 1 or 0 0 or 1 1* 1* (4) When Q = 1, Q = 0, the RS FF is “set”,
—
Race when Q = 0, Q = 1, ff is “reset”.
condition
(5) Will not change states unless opposite
More sophisticated Flip-flops. input triggered.
The D Flip-Flop (6) Only opposite state trigger causes state
change.
The DFF is bistable circuit with one input
plus the clock. The term “D” refers to “data”. (7) Race condition triggered by R = S = O
(S +)
(output cannot be predicted).
D S
Q Timing Diagrams
Clock It is important when designing sequential
R
Q circuits to understand the timing relationships
D
(R +)
between circuit elements.
Normally done by plotting the “transitions”—
Since the DFF is limit to one input, there is no the change between logic “1” and logic “0”—of
chance that a race condition will occur. the element.
The clocked DFF can be created simply by
replacing the “reset” input with an inverted “set” All the timing normally related to the clock.
input in the clocked RS Flip-Flop. In the fundamental timing diagram, unless
In this configuration, we can see that if the D specifically instructed to do otherwise, always
input = 1, when the clock is high, the FF will go start with the clock.
into the “set” state (“set” input = 1, “reset” = 0). Always assume a “50/50” clock, if duty
When the D input = 0, then the FF goes into cycle not specified.
the “reset” state when clock goes high (“set” input To illustrate DFF timing by plotting the DFF
= 0, “reset” = 1). “Q” output.
When clock = 0, the FF is Idle. Remember that a DFF output is normally
There is no “forbidden” state here-the DFF triggered when clock =1.
output is defined for All inputs. Assume DFF is originally “reset”, then an the
Inputs New Output rising edge of the clock, the DFF output goes to
— — “set” (1), when the D input is 1.
Clock D Q Q Q Q
On successive clock pulses, the DFF “Q”
0 X 1 or 0 1 or 0 Same Same changes as D changes.
1 1 0 1 1 0 Principle—Output “Q” of a simple DFF
1 1 1 0 Same Same tracks D when the clock ticks.
1 0 1 0 0 1
1 0 0 1 Same Same Output

So a flip-flop is useful as “memory”. It can Input

be set to 1 or 0 and retain that value until triggered 1
to the opposite state. Clock 0
470 | Physical Sciences

“Master-Slave”, or Delay Flip-Flops D Flip-Flop Symbols

It is often desirable to have a flip-flop whose
output does not change immediately when its
internal state is altered from “set” (Q = 1) to
“reset” (Q = 0) or vice-versa.
This sort of ff is called a “master-slave” or
“delay” ff.
The idea behind the master slaves ff is to have
a “master” (i.e., controlling) ff change states one
edge of a clock pulse (normally the leading edge)
and have a second ff connected to the first change
to the same state as the “master” on the trailing
edge, or “back side” of a clock pulse.
In this way, the internal state of the ff changes
one half clock cycle prior to the time in which
the changed state appears on the circuit
outputs.
D
Q if S = R = 0, output may be indeterminate.
Clock
Q Not useful to have a FF that will not have
Master Slave
DFF converted to master-slave type.
The slave (basically a clocked RSFF) always
mirrors the state of the master.
1 K R Q
The slave circuit changes state cycle after
2
the master. The device still operates as a D-FF;
no indeterminate state.
Slave FF(G)
Master FF
D
JK master slave FF designed as—
Flip-Flop details is not usually shown in
diagrams.
One symbol for a DFF is shown.
⇒ There is no small circle on the either
input.
∴ “1” is the active state (when D Q
clock and D = 1, output will → 1).
DFF’s with asynchronous Set C
and Reset are also available.
Circles on S and R inputs mean
that Set and Reset are “negative- D S Q
true” signals (active at level 0).
C R Q
⇒ “Q-not” out is also
available.
Set and Reset have the same problem, before
The JK Master-Slave FF
indeterminate outputs when S and R inputs are
both 1 simultaneously.
J S Q
C
The JKFF, sits those. The J input corresponds
to “set” while K corresponds to “Reset”.
The JK FF is designed so that the of J = K = 1
does not result in an indeterminate output when
clocked.
There may still be asynchronous RS inputs
with the usual cautions. However, the JK inputs
are not restricted.

OR Q
J

K
Q

Clock
The key states are J = K = 1, for either output state (set or reset).
If Q = 1 (“set”) and J = K = 1, output of the OR = 0, so the ff will be set.
(For J = 1 and K = 0 or J = 0 and K = 1) normal set of reset occurs.
Then for J = K = 1, when the clock ticks Q-the output state.
 Physical Sciences | 471

Input Outputs
AND OR New ↓
output output output JK +ve edge triggered
— —
J K Q Q 1 2 3 4 Q Q
0 0 0 1 0 0 0 1 Same Same J K C Outputs
—
0 0 1 0 0 1 1 0 Same Same Q Q0
0 1 0 1 0 0 0 1 Same Same —
0 0 ↑ Q0 Q0 No
0 1 1 0 0 0 0 1 0 1 change
1 0 0 1 1 0 1 0 1 0 0 1 ↑ 0 1 Reset
1 0 1 0 1 0 1 0 Same Same
1 0 ↑ 1 0 Set
1 1 0 1 1 0 1 0 1 0 —
1 1 1 0 0 0 0 1 0 1 1 1 ↑ Q0 Q0 toggle

Solved Questions
Q. 1. The load resistance of a full wave Q. 2. Explain the function of bridge
rectifier is 103 Ω . And the input AC voltage rectifier. Also write its advantage over single
(rms) is 200-0-200. Calculate average d.c. face full wave rectifier.
voltage, average d.c. current and ripple factor.
Solution : A bridge rectifier consists of four
Diode do not show any significant resistance
identical diodes arranged in a cyclic manner as
itself.
shown in the figure—
Solution : The peak value of d.c. current can
be calculated by P A
V
Im = m VPQ D1
RL D2 E
D D4
B
Vrms √
2 D3
= VRL
RL Q C
200 √2
=
103 F
= 0·2828 amp.
The average d.c. current is given by
2Im
Idc =
π
2 × 0·2828
=
3·14
= 0·18 amp
Average d.c. voltage
Vdc = Idc × RL
= 0·18 × 103
= 180 volt


√
2
Vrms
and Ripple factor, γ = 2 –1
Vdc


√ (200
180)
2
= –1

= 0·4843
 Unit-7

EXPERIMENTAL METHODS
Vacuum Technique Therefore, the rate at which pressure is
In experimental physics there is always a reduced is
need of vacuum of high degree. There are too dp (p – p0) dV
many techniques to create the vacuum. The =
dt V dt
electron microscope and Coolidge tube are same
device calling for the high vacuum.
High vacuum pumps may be roughly
= –s ( )
p – p0
V
classified into two categories : dp s
i.e., = – dt
(i) Those which remove the gas from the (p – p0) V
vessel to be evacuated and eject it directly into the If p1 and p2 denote the pressure in the vessel
atmosphere. These are usually mechanical pumps at time t1 and t2 respectively then by integration of
and are called backing pumps. One such pump both sides of the above equation
used commonly is rotary oil pump.
∫ ∫
p2 dp s t2
(ii) Those begin to operate below a certain – = dt
limiting pressure. In this case a fore-vacuum is p1 p – p0 V t1
required which can be produced by pump of first (p1 – p0 ) s
kind. or log = (t – t ) …(2)
(p2 – p0 ) V 2 1
Pumping Speed If p0 is negligible small, the above equation
The speed of pump at a given pressure is the reduced to
volume of gas pumped out from a vessel at that V p
pressure per unit time, the volume being measured S = log 1 …(3)
(t2 – t1 ) p2
at that pressure.
which gives us a definition of pumping speed.
Mathematically, If t1 – t2 = 1 sec
dV p1
S = …(1) and = e
dt p2
If p 0 is the lowest pressure attainable by a then, S = V
pump, p is the pressure at any instant t in the
vessel to be evacuated, V the volume of vessel, Hence, the speed of a pump is numerically
then applying Boyle’s law equal to the volume in which it can reduce the
1
(p – p0)V = (p – p0 – dp) (V + dV) instantaneous pressure to in one second. So its
e
neglecting the product of second order infinite- unit be litres/second.
simals
The Mechanical Pump
0 = (p – p0) dV – Vdp
The mechanical pump which is most
(p – p0) commonly used as a fore-pump is rotary oil
dp = dV pump.
V
524 | Physical Sciences
It consist of a solid steel cylinder called the
rotar which rotate centrically inside a thick walled
hollow cylinder, the stator.
A sliding vane V is always kept pressed
against the wall of rotar by the pressure of a rod
connected to a spring. The vane slides up and
down in a vertical groove and cuts off any
connection between inlet and outlet.
The vane divides the air space between the
stator and the rotar into two parts p and q. I is the
inlet port and E is the exhaust valve in the outlet
tube.
V
E I
q p
(a) (b)
V
I E I To
E
q p
p
(c) (d)
To work the pump, the vessel to be evacuated
is connected to the inlet port and the rotar driven
at high speed by an electric motor. If the rotor
rotates in the direction shown by the arrow, the
point of contact between it and the stator moves
clockwise, as shown in four successive positions
of the rotar (a), (b), (c) and (d).
When the rotar is in position (a), the space q
between stator and rotar is in communication with
the vessel and is filled with the gas (or air).
In the position b and c, the space q is more
and more diminished and air in it is compressed.
Simultaneously fresh air (gas) from the vessel fills
the space p behind the rotar. In the final position
(d), the pressure of air in q becomes greater than
the external pressure on E, E now opens and the
compressed air (or gas) is expelled. The process
continues until the pressure in the vessel becomes
too low, so that the compressed air (gas) is unable
to open the exhaust valve E. The vacuum attainable
is of the order of 10–2 mm of mercury.
The Diffusion Pump
If a face-vacuum of 0.11 applied this type of
vacuum pump can produce a pressure as low as
10–6 mm.
To fuse
pump
B
To vessel C
to be
exhausted
A
Mode of Working
As its name suggests pump works on the
principle of diffusion. A schematic diagram of the
diffusion pump is shown below through which we
can understand the action of pump.
vessel
Water
Liquid Water
air tap
To fuse
pump
Fig. : A schematic diagram showing the action of a
diffusion vacuum pump
A regular stream of air free mercury vapours
moves upwards along the vertical tube AB. To
vessel to be exhausted is connected to AB by a
side tube which is kept at a low temperature by
surrounding with a cooling jacket. The vapour
stream carries with it any air present in AB. Now
the concentration and partial pressure of air is
greater at D than at C since the mercury vapour is
air free and there is little or no air near C. Hence
air diffuses from the vessel into AB where it is
swept again by the vapour stream and removed by
the fore pump. The diffusion of air from the vessel
into AB continuous until the partial pressure on
AB and vessel are equal.
The mercury vapours also tends to diffuse
from AB where its partial pressure is higher

towards the end where its partial pressure is lower.
To minimise this tendency the fore pump lowers
the pressure at B. So that the mean free path of the
molecules there becomes greater and the orifice C
of the side tube is made narrow, its dimensions are
smaller than the mean free path of the air
molecules in AB. The diffusion velocity of
mercury vapours towards the vessel is much less
than the diffusion velocity of air from the vessel
towards AB. Thus, the possibility of back diffusion
of mercury is reduced. Any mercury vapours
proceeding through the side tube is condensed by
the cooling jacket before reaching the vessel.
Many design of this pump has been cons-
tructed and put to use. In general the mercury at
the bottom is heated so that the mercury vapours
rises upwards and is condensed around the cooling
jacket, the condensed mercury returns to the
bottom to be evaporated again.
Low Pressure (Vacuum) Measurements
The vacuum pressure gauges are used for
measurement of pressure below that of atmosphere
and this pressure is commonly referred as vacuum.
The science of low pressure measurement is
specialised field that required considerable care on
the part of the experimentalist two commonly
used units of vacuum measurement are the torr
and micrometer. One torr is pressure equivalent to
1 mm Hg at standard conditions and 1 micrometer
is 10–3 torr. The wide range of pressure to be
measured ranging from the normal atmospheric
pressure of 760 torr down to 10 × 10–9 torr makes
the problem of measurement very extensive.
There are two basic methods of measurement
of low pressure.
(i) Direct methods—The direct methods of
measurement involves measurement of a displace-
ment produced by elastic pressure transducers as a
result of application of pressure.
(ii) Indirect or inferential methods—
Manometer and Bellows gauges can be used for
pressure upto 0.1 torr, Bourdon gauges to 10 torr
and diaphragm gauges upto 10–3 torr. For the
measurement of pressure below these ranges other
types of gauges are used. These methods involve
the measurement of pressure through measure-
ment of certain other properties, which depends
upon the pressure to be measured. Some of the
properties are pressure dependent and are used for
Physical Sciences | 525
measurement of pressure are : change in volume
and change in thermal conductivity.
Pressure is measured by direct method using
spiral bourdon tubes, flat and corrugated diaph-
ragms capsules and may form of manometers.
They can be conveniently used for measurement
of pressures down to about 10 mm of mercury.
For measurement of pressure below this value
indirect methods have to be used. Indirect or
inferential methods for pressure measurement are
used in situation where it is possible to use a
phenomenon or property which is a function of
pressure. The phenomenon or property to be used
must be sensitive to changes in pressure in the
range selected. The commonly used devices for
the measurement of low pressure are McLeod,
Kundsen, viscosity, thermocouple, pirani and
ionization gauges.
Thermal Conductivity Gauges
As the pressure of a gas becomes smaller and
smaller, the mean free path, λ, of the molecules
becomes greater and greater. In fact at very low
pressures, the mean free path may be significantly
greater than the pertinent dimension of the
apparatus. Therefore, at these low pressures, a
linear relationship between thermal conductivity
and pressure is predicted by the kinetic theory of
gases. The pertinent dimension is the spacing
between two surfaces, one of them is hot and other
is cold.
Again when the pressure is increased
significantly, the thermal conductivity becomes
independent of the gas pressure. The transition
region, between dependence and non-dependence
of thermal conductivity on pressure is approxi-
mately in the range of 10 × 10–3 to 1 torr for an
instrument of a convenient and manageable size to
construct and use.
There are complication in the use and
operations of instruments which works on the
principle of variation of thermal conductivity with
pressure which leads to errors in measurement.
One complication arises from the fact that
radiation is another mode of heat transfer between
hot and cold surfaces which are simultaneously
present along the conduction. Most of the gauges
use a heated element which is supplied with a
constant energy input. This element reaches a
thermal equilibrium i.e., it reaches a steady
temperature when the heat energy supplied to it
526 | Physical Sciences
just equal to heat energy lost by it by conduction
and radiation. The conduction losses vary with the
gas composition and pressure. Thus for a given
gas the equilibrium (steady state) temperature of
heated element become a measure of the pressure
and in order to know the value of pressure, the
temperature must be measured. The temperature is
true indication of pressure in case the heat lost is
only on account of conduction. However, in
practice, heat is lost on account of radiation and
conduction both. In fact if the heat lost on account
of radiation is a major part of the total heat lost,
then the pressure induced conductivity changes
will cause a slight change in temperature and,
therefore, there will be very serious errors in the
measurement of pressure because it is assumed
that the heat lost is purely on account of
conduction. Therefore, in order to reduce this
error, heat losses due to radiation must be
minimised by using surfaces of low emissivity and
by making the temperature of cold surface as low
as practically possible. As heat losses on account
of both conduction and radiation depend on the
temperature of both hot and cold surfaces, the cold
surface must be maintained at a known constant
temperature if an overall accuracy in measurement
is to be achieved.
Another source of error is the heat conduction
loss through any solid supports on which the
heated element is mounted. The relative impor-
tance of above two losses varies with the details in
construction
The most commonly used types of conduc-
tivity gauges are of the following two types :
(a) Thermocouple gauge, and
(b) Pirani gauge.
(a) Thermocouple gauge : The basic elements
of a thermocouple vacuum gauge are shown in a
schematic form given below. The hot surface is a
thin metal strip whose temperature may be varied
by altering the value of current passing through it.
The temperature of the heated strip lies between
50°C to 400°C. The temperature of the metal strip
is measured with the help of a thermocouple which
is welded to the hot surface.
The thin metal strip and the thermocouple are
enclosed in a glass tube which is sealed into the
vacuum system whose pressure is to be measured.
The cold surface is formed by the glass tube
which is at the room temperature. Usually, the
accuracy of these gauges is not high enough which
may necessitate the measurement or correction for
changes in room temperature.
P
Cold surface
(glass tube)
Hot surface
Thermocouple
PMCC
Instrument
Fig. : Thermocouple vacuum gauge
The thermocouple provides an output voltage
which is a function of temperature of the heated
strip which in turn is a function of pressure of the
gas whose pressure is being measured. The output
voltage of the thermocouple is measured with the
help of a PMCC meter whose scale can be directly
calibrated in terms of pressure being measured.
Thermocouple gauges of one type or another
are available to measure pressure in the range
0.1 × 10–3 to 1 torr.
The advantages of the thermocouple gauge is
that these are inexpensive. However these are
subject to burn out if exposed to atmospheric
pressure when hot and cold currents are flowing.
(b) Pirani gauge : The pirani gauge like
thermocouple gauge is a device that measures the
Pi
Resistance
elements
Glass tube
To
bridge
circuit
Fig. : Pirani gauge
pressure through changes in thermal conductivity
of a gas which changes the temperature of a
heated element.

Unlike in the case of thermocouple gauge, the
function of heating and temperature measurement
are combined in a single element. The resistance
element is in the form of four coiled wires of
tungsten or platinum connected in parallel and
supported inside a glass tube to which the gas,
whose pressure is to be measured is admitted as
shown in the figure.
The cold surface is the glass tube. The pirani
gauge is connected in one of the arms of a Wheat-
stone bridge for the measurement of the resistance
of its heater element as shown in the figure given
below :
Measuring
element
To vacuum
space, P
Zero
setting
potentio-
meter Compensating
element
Fig. : Bridge circuit for Pirani gauge.
To identical pirani tubes generally connected
in the bridge circuit. One of the tube is evacuated
to a very low pressure and then sealed off while
the other has get admitted to it.
The evacuated tube acts as a compensator to
reduce the effect of bridge excitation voltage
changes and temperature changes on the output
reading.
In order to balance the bridge initially, the
pressure in the measuring element is made very
small and the balance potentiometer is used to set
the zero. Current flowing through the measuring
element heats it to a temperature depending upon
the gas pressure. The electrical resistance of the
element changes with temperature, and this
resistance change causes bridge unbalanced. The
magnitude of the output of the bridge produced on
account of the bridge unbalance is indicative of
the change in the value of resistance the heater
element and hence the unknown pressure.
The pirani gauges are rugged, inexpensive
and usually have a higher degree of accuracy
rather than thermocouple gauge. These can
measure pressure ranging from 10 × 10–6 to 1 torr.
On the other hand these requires calibration
against some pressure standards. Also the calibra-
Physical Sciences | 527
tion is non-linear and varies from one gas to
another. The transient response of pirani gauges
is poor. The time required for obtaining thermal
equilibrium may be of the order of several minutes
at low temperature.
Production of Very Low Temperature
For a long time it was considered that nothing
could be colder than ice and hence the temperature
of ice was chosen as the zero (starting point) of
the scale in all earlier systems of temperature
measurement. Later on, it was found that it is
possible to obtain the temperature well below the
ice point. The study of the laws of perfect gases
tells us that the lowest conceivable temperature on
the centigrade scale is 273 below the melting point
of ice. This is taken as zero of absolute scale or
absolute zero.
The production of very low temperature is
intimately linked with the liquification of gases;
because all experiments in the region of very low
temperatures involves the gases in liquid or in
solid form. For producing very low temperatures
the usual method is to be first liquify the suitable
gas by purely mechanical means and then by
evaporating it, the latent heat of vaporisation is
obtained at the expense of the heat content of the
liquid itself and intense cooling results.
The following are the methods employed to
produce very low temperatures :
(a) By adding salt to ice.
(b) Cooling by evaporation of a liquid under
reduced pressure.
(c) By the adiabatic expansion of a gas.
(d) Joule-Thomson effect.
(e) Regenerative cooling.
(f) By the process of adiabatic demagnetis-
ation of a paramagnetic salt.
The last method is useful only in the region of
liquid helium temperature and is used in producing
the lowest temperature.
Thermodynamics of Refrigeration
A refrigerator is a machine that operates in a
manner opposite to that of a heat engine. In it the
working substance abstract heat from the low
temperature source (The space to be refrigerated),
then some work is performed on it and finally it
rejects heat to a high temperature source (the
528 | Physical Sciences

surrounding atmosphere). Thus, refrigerator is a expenditure of some external work. We may call
device to transfer heat from a low temperature the heat extracted as output and the work as the
source to a high temperature source with the help input. It is convenient to express the efficiency of
of some external work. The working substance a refrigerator in terms of coefficient of per-
which undergoes a cyclic changes in the refrig- formance which is abbreviated as ‘Cop’ and is
erator is called the refrigerant and the process defined as the ratio of the amount of heat absorbed
itself is refrigeration. from the cold reservoir to the work done in
The schematic arrangement of a refrigerator is running the machinery.
shown below. The refrigerant is in the form of a i.e., the coefficient of performance,
stored liquid at a temperature equal to that of a Amount of heat absorbed from the cold reservoir
condenser whose temperature is equal to that of β=
Work done in running the machinery
the hot source. The refrigerant from the liquid If Q2 is the amount of heat absorbed from the
storage passes through throttling valve suffering
an adiabatic or Joule-Thomson expansion and cold reservoir and Q1 is the amount of heat
cools to a lower temperature then it passes through rejected to the hot reservoir during one complete
the evaporator and evaporates at constant tem- cycle, then the work done in running machinery,
perature and pressure absorbing required latent W = Q1 – Q 2
heat from the object to be cooled. so that the coefficient of performance,
Q
Hot reservoir β = 2
W
Q1
Q2
Condensor
=
Q1 – Q2
But in Carnot’s cycle
Q1 T
Liquid = 1
storage Compressor Q2 T2
where T2 is temperature of the cold reservoir and
Throttling T1 is the temperature of the hot reservoir.
valve
T2
Therefore, β =
Evaporator T1 – T2
If η is the efficiency of a Carnot’s engine
Q
operating between the same temperatures, we
Cold reservoir have
T2
Fig. : Schematic diagram of a refrigerator η = 1–
T1
These vapours then pass through the com-
T – T2
pressor where they are compressed adiabatically = 1
by a compression pump the heat of compression is T1
rejected at constant pressure by passing the 1 T1
or =
vapours through the condenser maintained at a η T1 – T2
low temperature, so that the working substance of
refrigerant leaves the condensor as saturated or 1 T1
or –1 = –1
sub-cooled liquid and the fresh cycle of operation η T1 – T2
starts. 1–η T – T1 + T2
or = 1
Coefficient of Performance of a η T1 – T2
Refrigerator T2
=
In a refrigerator it is desired to extract a large T1 – T2
amount of heat from a cold body with the = β

The coefficient of performance
1–η
β =
η
Just as no heat engine can be more efficient
than Carnot’s heat engine, the ‘cop’ of no
refrigerator working between two reservoir can be
higher than that of Carnot’s refrigerator similarly.
Coefficient of performance of actual refrig-
erators vary from about 2 to about 6.
Adiabatic Demagnetisation of a Para-
magnetic Salt
In 1926, Debye and Giauque should theoret-
ically that the temperature considerable below 1 K
could be obtained by the process known as
adiabatic demagnetisation of a paramagnetic salt’
(i.e., those substances for which the magnetic sus-
ceptibility χ is small; but positive). The principle
of method is as follows :
The process of magnetising a substance
involves doing work on the substance in aligning
the elementary magnets in the direction of
magnetic field. This work is stored in the substance
in the form of magnetic energy. If the magnetised
substance is demagnetised adiabatically, it has to
do work and the required energy for it is drawn
from within itself. As a result the substance cools,
considerable cooling can be produced by employ-
ing strong magnetic field and low initial
temperature. This follows from Curie’s law,
according to which the magnetic susceptibility
of a substance is inversely proportional to the
1
absolute temperature i.e., χ ∝ or χ = CT. The
T magnetisation per gram mole I of the system
D1
D2
A work on the substance in aligning the atomic
N S
Paramagnetic salt
Liquid helium
Fig. : Schematic diagram of cooling through
adiabatic demagnetisation of paramagnetic salt
method.
final temperature is determined by measuring
the magnetic susceptibility of the substance and
then calculating from Curie’s law method.
Physical Sciences | 529
Debye, Giauque and Macdougall were able to
produce the temperature below 1 K with the help
of paramagnetic salt godolinium sulphate. The
experimental arrangement is shown in the adjacent
figure. The sample (Paramagnetic salt godolinium
sulphate, say) is suspended in a vessel surrounded
by liquid helium contained in a Dewar flask D1
which in turn is surrounded by a Dewar flask D2
containing liquid hydrogen. The whole arrange-
ment is placed between the poles of a strong
electromagnet.
Now the following procedure is adopted—
The magnetic field is switched on, so that the
speciman is magnetised. The heat due to mag-
netisation is removed by first introducing
hydrogen gas into A and then pumping it off with
a high vacuum pump, so that the speciman is
thermally isolated. In the meanwhile the speciman
picks up the temperature of the liquid helium i.e.,
the speciman and the liquid helium are at same
temperature. Then the magnetic field is switched
off. Adiabatic demagnetisation of the speciman
takes place and its temperature falls. The tempera-
ture of the speciman is determined by fitting a co-
axial solenoid coil round the tube A and measur-
ing the self inductance and hence susceptibility of
the substance at the beginning and at the end of
experiment.
Hass in 1941, was able to produce the
temperature upto 0.002 K using a double sulphate
of potassium and aluminium theory.
When a paramagnetic salt is placed in a
magnetic field of strength H, the intensity of
change by dI. During the process of demag-
netisation, the external magnetic field has to do
magnets in the direction of the field. The work
done by the field on the substance.
= HdI
Therefore, the work done by the substance
= –HdI
According to first law of thermodynamics
dQ = dU + dW …(1)
where d Q is the heat supplied to the system which
results a change dU in its internal energy.
Here, dW = PdV – HdI
Therefore, equation (1) gives
dQ = dU + PdV – HdI …(2)
530 | Physical Sciences

But change in volume of a solid for small Substituting this in equation (9), we get
change in temperature is negligible.
i.e. d = 0
then equation (2) became
(∂V
∂H) T
= T (∂T∂I ) H
+H (∂H∂I ) T
…(11)

dQ = dU – HdI …(3) Using above equation, we can calculate the

internal energy or the entropy in an arbitrary field
According to second law of thermodynamics
from the internal energy or entropy in zero field at
dQ = TdS …(4) the same temperature, when the magnetisation
where d S is the change in entropy. curves at different temperatures are known. From
Using equation (3) and (4), we have equation (3) and (6), we have
TdS = dU –HdI
or dU = TdS + HdI …(5) dQ ={(∂U
∂T )
∂V
d T ( ) dH} – H
∂H
H T
In these equations the entropy S, the intensity
of magnetisation I and internal energy U are to be
regarded as function of absolute temperature T
{(∂H∂I ) dH + (∂T∂I ) dT} T H
and the magnetic field strength, H thus we have ∂U ∂I ∂I
= [( ) dT + { T( ) + H( ) }dH]
dS = ( )∂S
∂T H
dT +
∂S
( )
∂H T
dH  ∂T H ∂T
∂I
∂H
∂I
H T

– H{ ( ) dH + ( ) dT}
dI = ( ) ∂I
dT + ( )∂I
dH  ∂H ∂T T H


∂T H ∂H T …(6)
∂V ∂V [Using (11)]
and dU = ( )∂T H
dT + ( )
∂H T
dH
{(∂U ∂I
∂T )
∴ dQ = – H ( ) }dT
Substituting values of dS, dI and d V in H∂T H

equation (5), we get ∂I

+ T ( ) dH …(12)
∂V ∂V ∂T
( )
∂T H ( )
dT +
∂H T
dH
Now the thermal capacity of the substance at
H

∂S ∂S
=T {( ) ( ) }
∂T H
dT +
∂H T
dH
constant field may be expressed from equation
(12) to be
∂I ∂I
{( ) ( ) }
+H
∂T H
dT +
∂H T
dH …(7) CH = (∂Q
∂T ) H
∂U
=( )
∂T H
–H (∂T∂I ) H
…(13)
Comparing coefficients of d T and dH on both Equation (12) can be written as
sides, we get

( )∂U
∂T H
= T ( )∂S
∂T H
+H ( )∂I
∂T H
…(8) dQ = C HdT + T (∂T∂I ) H
dH …(14)

∂V ∂S ∂I During adiabatic demagnetisation no heat

( )∂H T
= T ( )∂H T
+H ( )∂H T
…(9) enters or leaves the system i.e., dQ = 0, so that
equation (14) gives
As U is the function of the state of the
system, dU must be a perfect differential, i.e.
∂ ∂V ∂ ∂U
CH d T + T (∂T∂I ) H
dH = 0 …(15)
[ ( )] [ ( )]
∂H ∂T H T
=
∂T ∂H T H According to Curie’s law
∂ ∂S ∂I
or[ { ( ) ( ) }]
∂H
T
∂T H
+H
∂T H χ
C
T
= …(16)
∂ ∂S ∂T
[ { ( ) ( ) }]
=
∂T
T
∂H T
+H
∂H T
where C is a constant known as Curie constant
and χ is molecular susceptibility.
∂I ∂S
( ) ( )
∂T H
=
∂H T
…(10) But χ =
M IV
=
H H
 Physical Sciences | 531

where V is volume occupied by 1 gm mol of the entropy diagram is illustrated. Here the entropy of
substance and M is its molar magnetisation, this the substances is represented schematically as a
gives, function of temperature for different external
fields H = O and H = H. The change in entropy
H
I = .χ during isothermal magnetisation is represented by
V the arrow AB. This shows that during magnetisa-
C H tion the entropy of the substance decreases. Now
= .
T V if the field is switched off, the substance demag-
C H netises adiabatically, so that the temperature of
= . substance falls while entropy remains constant as
V T
represented by arrow BC.
kH
= …(17)
T S
A H=O
C
where k = is the Curie constant/unit volume
V
differentiating equation (17) w.r.t. T at constant
field H, we get C
B H=H
∂I
( )
∂T H T
k
= – 2H …(18)
T
Substituting this in equation (15) we get 1K

C H dT + T (– Tk ) 2 H
dH = 0 With the help of this process, temperature of
the order of 10 –4 can be reached. The measure-
ment of temperature in this region is a matter of
kH
or TdT
dH = great difficulty but what is actually done is to
CH assume the validity of Curie’s law and calculate
Integrating above equation between initial from the measurement of the susceptibility of the
and final, states, we get, C
sample at Curie temperature T* = (also called as
T2 0 χ
∫ T1 TdT ∫ H CH dH
kH
= magnetic temperature). Later from Curie tempera-
ture, the thermodynamical temperature T is
2 2 0
or [T2 ] T2

T1
=
k
CH [H2 ] H
determined. Steerland; Gorten and de-Klark were
able to reach a magnetic temperature T* = 0.033K.
While working with powdered mixed crystals of
Here, we have replaced CH within the integral chromium alum and aluminium alum which
– according to oxford workers corresponds to
by its average value C H
0.010K. It may be noted that attainment of
2 2 k
T2 – T1 = – – H2 absolute zero is an impossibility.
CH
2 2 k
Measurement of Low Temperature
or T2 = T1 – – H2 The property used to measure temperature
CH must be measurable with good accuracy with
 k H 2  1/2 relative ease. It should be sensitive, i.e., it should
or T2 = T1  1 – –

CH
T1 () …(19) exhibit relatively large percentage change with
change in temperature over the required
We see that T2 > T1 i.e., the temperature of the measurement range.
salt falls during the process of adiabatic demag- Vapour pressure thermometry where basis is
netisation. The fall in temperature is greater for dP L
the Clausius. Clapeyron equation = can
large values of the magnetising field H and low dt T∆V
values of initial temperature. In the figure given be used to liquid helium temperature, but its use is
below the cooling process using a temperature very time consuming.
532 | Physical Sciences

Thermocouple can be used from room (i) The pressure bulb contains the thermometer
temperatures to a few degree Kelvin. Calibration fluid and is placed where the temperature to be
equations of the type measured. It attains temperature equilibrium with
E = a0 + a1 , ∆t + a 2 (∆t) 2
where ∆t is the temperature difference between
the two junctions gives the thermo emf. Thermo-
couple junction can be made very small and have
a fast response time and minimum thermal load.
Their output voltage is readily measured by

100
modern high impedence digital voltmeters and
digital outputs may be fed into a computer which Receiving
can be programmed with the calibration equation element
0 (bourdon tube)
to give a direct temperature read out. The
coefficient can be determined using a few fixed -100
calibration. Calibrated
scale
Thermocouple voltage decreases rapidly at
low temperatures. But Au Fe-chromal (Au com- Pressure
ponent (0.02 – 0.07%) gives acceptable perfor- bulb
mance down to about 10µV/k at 4°K.
Capillary
Carbon resistance thermometers are useful in tube
the liquid helium range and down to 1°K. They
require a maximum of these calibration points. Fig. : Generalised filled system diagram
They have rather low thermal conductivity and so its surroundings, thereby developing a given
care must be taken to arrive self-heating effects pressure or displacement of liquid.
and allow the resistor time to reach thermal
(ii) A capillary tube is used to connect the
equilibrium.
bulb to a read out or to a transmitting system.
Pressure Gauge Thermometers (iii) The receiving element is a device that
The pressure system thermometers are useful respond to changes in pressure or volumetric
in the temperature measurement in the range – 270 expansion of liquid caused by changes in tem-
to 760°C with a 0.5 – 2.0% accuracy. perature.

These thermometers essentially consist of a Transducers

bulb, containing a liquid, gas or vapours, which
is immersed in environments. The bulb is
connected by means of capillary tube to some
pressure measuring device, such as Bourdon tube.
An increase in temperature cause the liquid or
gas to expand, thereby increasing the pressure on
pressure measuring device. This pressure is thus
taken as an indication of temperature.
Since the operation of pressure gauge ther-
mometers is based upon the expansion of liquid,
gas or vapour filled in the bulb of thermometers,
these temperature measuring systems are com-
monly known as filled system. The generalised
configuration of all filled systems are identical as
shown in the figure ahead :
The filled system consists of the following
parts :
These are devices for converting a non-
electrical signal into electrical signal (or vice-
versa), the variations in the electrical signal being
a function of the input. Transducers are used as
measuring instruments and in the electroacoustic
field, the term being applied to gramophone pick
ups, microphones and loudspeakers. The different
electrical phenomena employed in transduction
elements of transducers are listed below—
(1) Resistive
(2) Capacitive
(3) Inductive
(4) Electromagnetic
(5) Piezo-electric
(6) Ionization
(7) Photo electric
 Physical Sciences | 533

(8) Photo conductive (ii) Poor maintenance.

(9) Photo-voltaic (iii) Change in process parameters irregu-
(10) Thermo electric larities, upsets etc.
(11) Potentiometric (iv) Poor design.
(12) Frequency generating (v) Design limitations.
(vi) Due to person operating the instrument.
Error Analysis
Normal or Gaussian Curve of Errors
In statistical terminology error are classified
as— Gaussian law of errors is the basis for the
1. Mean error— major part of study of random effects. This type of
n distribution is most frequently met in practice.
Σ | xi – –x | The law of probability states the normal
i=1
occurrence of deviations from average value of an
n
infinite number of measurements or observations
2. Standard error— can be expressed by
 n 1/2 λ
 Σ (xi – –x)2
 y = exp ( – h2 x2) …(1)
=  π
√
i=1
 n 
Note—The reader here is cautioned not to
Which is equivalent to error’s standard confuse x with magnitude of a quantity. Here x
deviation. means deviation.
3. Probable error—The probability of an
error falling within the limits µ – λ and µ + λ is where x = magnitude of deviation from mean.
exactly equal to the chance of an error falling y = number of readings at any deviation x
outside these limits, i.e., the chance of an error (the probability of occurrence of deviation x)
1
lying within µ – λ and µ + λ is . h = a constant called precision index.
2
2 Equation (1) leads to curve of type shown in
Probable error λ = 0·6745σ ~ σ the figure and this curve showing y plotted against
3
x is called ‘‘Normal or Gaussian probability
Types of Error are curve’’.
1. Gross errors—Which are due to human Another more convenient form of equation
mistakes in reading or in using instruments or describing Gaussian curve uses standard deviation
errors in recording observations. σ and is given by
2. Systematic errors—Which occurs due to 2
shortcomings of the instrument such as defective
or worn parts or ageing effects of the environment
y=
1
σ
√2π
exp – (2σx ) 2 …(2)

on the instrument and further classified as— h

(i) Instrumental errors y=—
y π
(ii) Environmental errors
(iii) Observational errors
3. Random errors—Which are due to 1 2 2 h
y =  e–x /2σ
2 2
unknown causes and are normally small and y = — e–h x
σ 2π π
follow the law of probability. These are treated
mathematically.
Sources of Errors
(i) Insufficient knowledge of process para- –∞ +∞
meters and design conditions. –3σ –2σ –σ –r 0 +r σ 2σ 3σ
534 | Physical Sciences
x2
∫ x1 exp (–h2 x2)dx
The second form of equation is particularly h
useful because σ is the usually known quantity of =
π
√
interest. In the figure the deviations from the mean
value are divided in terms of σ units, so that the The fraction of the total number of deviation
deviations are x = σ, 2σ, 3σ……etc. The Gaussian falling between zero and x are
curve represents the error distribution of any set of
x2
∫0
data obeying the normal law and may be used as h
no – x = exp (–h2 x2)dx
such once the value of σ has been determined. π
√
This curve is symmetrical about the arith-
metic mean value and area under the curve is Graphical Analysis and Curve Fitting
unity. Under the conditions specified here the total Significant information can be obtained from
number of readings taken is represented by 1. proper plots of experimental data. However, the
This can be explained as follows : Suppose for plotting of experimental data in correct fashion
h involves the understanding of physical phe-
the time being that we consider to be unknown nomenon involved.
π
√
and replace it by the symbol A. If we have a large Assuming that one knows what one wants to
number of readings n, the probable number ∆n, examine with the help of plotted experimental
with deviation between x and x + ∆x is given by; data. The graphical plots may be carefully
prepared and checked against the appropriate
∆n = ny ∆x theoretical results.
= n A exp (–h2x2) ∆x It is frequently desired to have a correlation of
the experimental data in terms of analytical
If we integrate the above expression for –∞ to variables that were measured in an experiments. It
+ ∞ we shall have all the cases, is very easy to obtain an analytical relation if the
+∞ data may be approximated by a straight line.
nA ∫ –∞ exp (– h2x2) dx = n However, it is not always easy to present the data
in the form of straight line as many other
+∞
∫ –∞ exp (– h2 x2) dx =
functional variations may be involved. However
Thus A 1 affords should be made that some form of
+∞
relationship is built up between variables so that
∫ –∞
h a linear graph be obtained. Linear graphs have
or exp (–h2x2) dx = 1
π
√ their obvious advantages and the most elementary
one is that it very easy to evaluate the results
Thus, the integral of y from – ∞ to + ∞ is equal from a straight line rather than from any other
to unity. The fraction of the total number of form.
readings occurring between the values x1 and x2
will be equal to the area under the curve between In actual practice, the experimental data may
these values of x. be in the form of a polynomial, exponential or any
other complicated logarithmic function. It is most
∫
h x2 convenient, however, to plot the data in such a
n1 – 2 = exp. (– h2x2)dx
π x
√ form that a straight line is obtained by applying
the appropriate functional relationship.
where n1 – 2 = number of readings occurring bet-
ween x1 and x 2 . If the area between x1 and x 2 is Several different types of functions and their
0.5, then 50 per cent of the deviations fall between plotting methods that are used to produce straight
x1 and x2 . lines on graph paper are given in the table drawn
below. The method of least square should be used
In general the probability for finding a devia- for graphical portrayal of the result in the form of
tion in an interval between x 1 and x2 becomes straight line. The various constants can be
x2 evaluated with the help of there straight line
P (x1 , x2) = ∫x ydx graphs.
 Physical Sciences | 535

Table : Methods of plotting various functions to obtain straight line

Function relationship Method of plot Graphical representation of

parameters
y
Slope a
y = ax + b y versus x on linear paper
b
x
O

log y
Slope = b
y = ax b log y versus log x on log-log paper
log a
x=0 log x
or log x = 1
log y
Slope = b log e
y = ae bx log y versus x on semi-log paper
log a

O x
1/y
x 1 1 Slope a
y= versus on linear scale Extra
a + bx y x
plot
b
O 1/x
y y1
y = a + bx + cx2 y – y1 x x1
versus x on linear paper Slope c
x – x1
b + cx1
x
O
x x1
y y1 2
y=
x
+C
x – x1 Slope = b + ab x1
a + bx versus x on linear paper
y – y1
a + bx1
x
O
y
[( (
log y
1
1/(x x1) [
1/(x – x1)
2
y = ae bx + cx log [(yy )
1
]versus x on semi- Slope = c log e

log paper b + cx1 log e

x
O
536 | Physical Sciences


√
Method of Least Squares
n
Suppose we have a set of n readings x1, x 2 , x 3 , and sb = s
n Σ xi2 – (Σ xi)2 y
……, xn the sum of square of their deviations


√
1
from some mean value xm is where sy = Σ (ax + b – y )
i i
2
n n
S = ∑ (xi – xm)2 Standard deviation of x is
i=1


√
Supposing we wish to minimise S with respect 2
to mean value xm we have,
n
sx = (y a– b – x)
1
n
Σ i
i

δs
δxm (
= 2 ∑ xi – n xm
i=1
) =
sy
a
= 0 6
n
1
or xm = ∑
n i=1 i
x
5
n y = 0.672x + 0.296
1
But ∑ xi = x
—
4
n i=1
∴ x = –x
m
3
y
Thus, the mean value which minimise the sum
2
of the squares of the deviations is the arithmetic
mean. This is a simple example of application of
1
the method of least squares. Another method of
least squares is given below.
0
Suppose that two variables x and y are 1 2 3 4 5 6 7
measured over a range of values. We want to x
obtain a simple analytical expression for y as a Fig. : Method of least square for curve fitting
function of x. The simplest type of function is a
linear one and hence we try to find the best linear Example—The iron losses in ferromagnetic
function connecting y and x. material used in a transformer vary with frequency.
Suppose the linear function is y = ax + b f, of the supply given to the transformer. For a
particular transformer these iron losses were
Now, we want to minimize the quantity, determined at various frequencies with a constant
n
flux density in the ferromagnetic material. The
S = ∑ [yi – (axi + b)]2 results are :
i=1
This is done by taking derivatives of S with Frequency, f in Hz 1100 1400 1700 2000
respect to a and b and setting them equal to zero. Iron losses, P, in mW 46 62 94 122
This gives
Assuming the iron losses to have general
nb + a Σxi = Σyi
form P = Af2 + Bf determine the constants A and
and b Σ xi + aΣxi2 = Σxi yi B to achieve best fit for the four measured values,
Solving the above equations, we get, using method of least requires.
n Σxiyi – (Σx i) (Σy i) Now P = Af2 +Bf
a = …(1)
n Σxi2 – (Σx i)2 or P/f = Af + B
(Σy i) (Σx i2 ) – (Σxiyi) (Σx i) Therefore, the above relationship can be
b = …(2) expressed as
n Σxi2 – (Σx i)2
The standard deviations of a and b may be y = ax + b
found as : which is a straight line relationship between
P


√ n y =
sa = s f
n Σ xi2 – (Σ xi)2 y and x = f

ATOMIC AND MOLECULAR PHYSICS
A free atom is ordinarily in its normal state of
lowest energy. In order that an atom may emit
radiation, it must, somehow or other, be sent to a
state of higher energy. This happens when an
atom is placed in a sufficiently strong electric
field or is hit by highly energetic electrons or is
made to collide rapidly with other atoms. When
this atom known as excited atoms falls back into
its normal orbit, a transition of definite amount of
energy takes place and it is accompanied by a
definite spectral line. The intensity of the spectral
line depends on the probability of the transition
between the states.
Spectra may be divided into sub-classes as
follows :
Spectra
| employed and for separation by refraction, prism
↓ ↓
Emission Absorption
| Both methods may be used except in the
↓ ↓
Continuous Discontinuous
| be diffraction of light, but on the other hand
↓ ↓
Line spectra Band spectra
Line Spectra
They contains discrete wavelengths distri-
buted throughout the spectra and the wavelengths
of the lines emitted are the characteristics of the
element concerned. This spectra is also known as
atomic spectra. An atom always emits its charac-
teristic radiation independent of atom’s form.
Band Spectra
The band spectra can be obtained when the
emitting substance is in the molecular state.
Therefore, also known as molecular spectra.
When the substance (molecular form) is
heated upto dissociation temperature, the bands
disappears because molecules dissociate them-
selves into atoms.
Unit-8
Each band contains several lines separated
from one another. The lines are wider at red end
and the spacing decreases at violet end.
Absorption Spectra
When the light from a source of continuous
spectrum is passed through an absorbing layer,
certain radiations are found missing which
appears as lines or bands in emergent light.
Fraunhofer lines in the continuous spectra of
sun is an example of missing wavelength or
absorption spectra.
Investigation of Spectra
The separation of component of light can be
established either by refraction or by diffraction.
For the separation by diffraction, gratings are
are used.
region below 1250 Å where a grating is necessary.
Furthermore both the methods depends upon the
wavelength, greater is the wavelength, greater will
greater is the wavelength smaller will be the
refraction of light.
The prism method has the advantage of
greater light intensity, whereas grating method
generally affects greater resolving power.
In far infrared region (FIR) thermopiles and
bolometers are used. However below 13,000 Å
photographic plates are used.
The Bohr Atom
The first successful explanation of atomic
puzzle was put forward by Neils Bohr in 1913.
The concept of matter waves leads in a natural
way to this theory, as de-Broglie found and this
was the route followed. Bohr himself used a
different approach, since de-Broglie’s work came
a decade later, which makes his achievements all
the more remarkable. The results are exactly the
same.
560 | Physical Sciences

We start by the behaviour of an electron Photon energy = Initial energy – Final energy
orbiting around the nucleus. The corresponding h ν = Ei – Ef
de-Broglie wavelength will be skiles
h
λ =
mv
…(1)
E1( 1 1
–
ni 2 nf 2 ) = Ei – Ef

Also the necessary centripetal force required

for circular motion is balanced by electrostatic
or Ei – Ef = – E 1 (n1 – n1 )
f
2
i
2

force between nucleus and electron. Ei – Ef

or ν =
mv 2 1 e2 h
= E
r 4πε0 r2
e
= – 1 (
1 1
–
h nf 2 ni 2 )
or v = …(2)

√ 4πε0 mr Since λ =
c
v
From (1) and (2)
E
1
= – 1 (n1 – n1 )

√
…(5)
h 4πε0 r λ ch n n
λ = f i
e m This equation states that the radiation emitted
Also according to Bohr the condition of by excited H atom should contain certain wave-
stability is lengths only.
n λ = 2πrn (n = 1, 2, 3, …) These wavelengths fall into definite sequence
depending upon the final quantum number n f.
or
nh
e 
√ 4πε0 rn
m
= 2πrn Since ni > nf in each case, we can categories the
following spectral series :
or the orbital radii in Bohr atom can be given by E

rn =
n 2 h 2 ε0
(n = 1, 2, 3, …) …(3)
1
Lyman nf = 1 = – 1
λ ( ) 1 1
ch 12 n2
– n = 2, 3, 4, …
πme2 E
The radius of the innermost Bohr orbit can be
1
Balmer nf = 2 = – 1
λ ( ) 1 1
ch 22 n2
– n = 3, 4, 5, …
calculated from (3)
E
ch (3 n )
1 1 1 1
a0 = r1 = 5·292 × 10–11 m Paschen n = 3 = –
λ
f – n = 4, 5, 6, …
2 2

and rn = n2 a0 [from (3)] E

ch (4 n )
1 1 1 1
Brackett n = 4 = – –
Again the electron energy En is given in λ
f 2 2
terms of the orbit radius rn as E
ch (5 n )
1 1 1 1
Pfund n = 5 = – – n = 5, 6, 7 …
Putting the value of rn . f
λ 2 2

me n E
En = –
1
()
8ε02 h2 n2
= 21
n
…(4)

(n = 1, 2, 3, …)
e2
En =
8πε0 rn
E1 can be calculated from equation (4)
E1 = –13·6 eV
E1 represent the energy carries by an electron
in Bohr’s first orbit. When an electron jumps from
one energy level to another, with the difference in
energy between the levels being given off all at
once in a photon rather than in some more gradual
manner, fits in well with the Bohr model.

The radiations are found in the regions as
follows—
Lyman → UV
Balmer → Visible
Paschen → IR
Brackett → IR
Pfund → IR
Shortcomings of Bohr’s Theory
(i) It could not explain the fine structure of
spectral lines of hydrogen atom.
(ii) It cannot explain the variation of intensity
of spectral lines of an element.
(iii) The theory is applicable to only one
electron atoms such as hydrogen, singly ionized
helium etc.
Spectra of One Electron Atom
The spectra of alkali metals, e.g., lithium,
potassium, sodium, rubidium and cesium is
identical to the spectrum of hydrogen atom in the
sense that all of these contain a single valence
electron. In hydrogen atom the energy of the
moving electron is expressed in terms of
R
Rydberg’s constant as – 2 whereas in case of
n The quantum number J has the possible values
R J = | L + S |, | L + S1 |, … | L – S |
alkali spectra this term value changes to T = – 2
n0
where n0 is known as effective quantum number.
The prime reason behind this change lies in the
penetration of valence electron in the atomic
trunk. The value of e = l ⇒ n – 1 the path of
electron is circle but for small values the chance
of penetration of atomic trunk, increases,
justifying again the above aspect.
Spectral Terms
The valence electron of an atom is respon-
sible for the optical spectrum of an atom. The
quantum state of an electron and that of an atom
are depicted by different values of l, s, j and L, S,
J respectively. Needless to say for one valence
electron atom the values of l, s and j coincides
with the values of L, S and J because the inner
shell electrons does not contribute to the optical
spectra due to their non-contribution to angular
momentum of atom. For many electron systems
the values of L, S and J are the vector sum of l, s
and j for different free electrons. The following
Physical Sciences | 561
notation is used for the spectral terms correspond-
ing to various l values :
l 0 1 2 3 4
Term s p d f g
Similarly the capital letters S, P, D and F
depict the state of an atom for the values of L as
0, 1, 2, … etc.
Total Angular Momentum (The Vector
Model)—In Classical Mechanics the total angular
momentum (orbital + spin) is an important quantity
because its rate of change is equal to the net torque
applied to the system. Similarly in wave mechanics,
the total angular momentum J found from the vector
addition
J = L+S
Plays an important role. Because the vector
model applies to many electron, as well as one
electron, atoms, we here introduce the following
notations—
Quantum numbers specifying states of individual
electrons are denoted by small letters. Quantum
numbers representing atomic states will be denoted
by capital letters. In the special case of a one
electron atom, the electronic state is the atomic state
and capital letters are used.
In this notation, then the magnitude of J is
quantized according to
|J| = √
J(J + 1) h
where L and S are orbital and spin quantum numbers.
As in the case for the orbital and spin angular
momenta, the component of J in a physically defined
z-direction is separately quantized. We have
J z = MJ h
MJ = J, J – 1, J – 2, …, – J
For a hydrogen like atom
J = { L + S‚ L – S L > 0
S L=0
In case of one electron system the value of S
1
is + , so its multiplicity r = 2S + 1 will be 2. So
2
single electron system always gives rise to a
doublet for the values L +( ) ( )1
2
and L –
1
2
for
J. But the ground state is always a singlet. In two
electron system S = 0, 1. Thus, the state is either
1
singlet or triplet, for three electron system S = or
2
3
, hence the state is doublet or quartet and so on.
2
562 | Physical Sciences
The atomic state of an atom can be
represented in terms of multiplicity r and total
angular momentum J. For example term 2P 3/2
3 the line for which l increases and j decreases or
represents r = 2, L = 1 (P state) and J = , so S
2 vice-versa is not allowed.
1 (II) Quantitative Rules—
will be + . Sometimes principle quantum number
2
n can also be included, e.g., for n = 2.
2 2 P 3/2 .
In one electron system every state is doublet
and may be written as 2 P 1/2, 2 P 3/2, 2 D3/2 and 2 D5/2
and so on. The ground state of a single electron
system is always singlet. The separation for 2 P, 2 D
and 2 F doublets are shown :
The separation among energy levels :
(i) Increases with increasing atomic number z.
(ii) Decrease with increasing value of n.
(iii) Decreases with increasing l.
Selection Rules
All the transitions among energy levels are
not possible. Certain selection rules (sometimes
known as transition rules) controls the transition
occurrence.
n can change by any value
l can change by ±1 only
J can change by 0 or ±1 only
S can change by 0 only.
In the presence of magnetic fields, for weak
magnetic field the change in mj
∆ml = 0 or ±1
For strong magnetic field, the coupling
between l and s is broken up,
∆ml = ±1
∆ms = 0
There are some intensity rules also which
governs the intensity of spectral lines. These can
be sub-divided in following two categories.
(I) Qualitative Rules—
(i) For brightest line j and l changes by same
value.
(ii) If there are more than one line for which j
and l changes by same value than the line
involving the highest j value will be strongest.
(iii) The line for which l decreases is stronger
than the line for which l increases.
(iv) l and j should change in either way, i.e.,
(i) The sum of intensities of all lines corre-
sponding to transitions from a given level is
proportional to the statistical weight (quantum
weight) of that level.
(ii) The sum of intensities of lines which end
on a common level is proportional to the quantum
weight of that level.
The quantum weight of a level can be given
by (2j + 1).
Coupling Schemes
Since an atom is consist of large number of
electrons with different orbital and spin momenta,
coupling scheme is necessary to obtain the
resultant orbital and spin momenta of an atom as a
whole-coupling scheme are of two types :
(a) L-S coupling (Russell Saunders coup-
ling)—In L–S coupling scheme several l vectors
combine to form a resultant L and S vectors a
separate resultant S. The resultant L and S then
combine to form a resultant J, which represents
the total angular momentum of the whole electron
system of the atom, i.e.,
(l1 + l2 + l3 + …) + (s1 + s2 + s3 + …) = L + S = J
In forming the resultant J certain quantum
rules must be satisfied :
(i) All the three vectors L, S and J must be
quantised.
(ii) L may have the values 0, 1, 2, 3,
… depending upon the number of electrons in the
atom and the direction of their orbital vectors.
1 3 5
(iii) S may assume values 0, , 1, , 2, , …
2 2 2
depending upon the number of electrons in the
atom and the directions of their spin vectors.
(b) j-j coupling—In this type of coupling,
each electron contributes to the total angular
momentum of the atom by combining first its
individual spin and orbit vectors by relation
j = l + s. The vector sum of each j corresponding
to each electron gives the total angular momentum
J of the atom, i.e.,
(l1 + s1 ) + (l2 + s2 ) + … = (j1 + j2 + j3 + …) = J

The diagrams for L-S and j-j coupling are
shown below :
Fine Structure of Hydrogen Lines
According to Bohr, the lines in the hydrogen
spectrum should each have a well defined wave-
length. Spectrograph of high resolving power
showed that the Hα, Hβ and Hγ lines in the hydro-
gen spectrum are not single. Each spectral line
actually consisted of several very closed lines
packed together.
Michelson found that under high resolution,
the Hα line is resolved into two close components
with a wavelength separation of 0·13 Å. This is
called the fine structure of spectral lines which
cannot be explained on the basis of Bohr’s theory.
Spin Orbit Coupling and Splitting of
Quantum State
The fine structure doubling of spectral lines
may be explained on the basis of magnetic inter-
action between the spin and orbital angular
moments of atomic electrons. This spin orbit
coupling can be understood in terms of a straight
forward classical model. In its orbital motion, the
electron goes in a closed path round the nucleus of
charge Ze. When viewed from the electron. It is
equivalent to the nucleus of charge Ze going
round the electron in a close path (see figure).
r v r
–e +Ze
(a) (b)
At the point where the electron is situated, a
magnetic induction is produced because of the
nucleus go round in a closed path. The interaction
between the electron’s spin magnetic moment and
Physical Sciences | 563
this magnetic field leads to the phenomenon of
spin-orbit coupling. The magnetic potential
energy associated with the spin-orbit interaction is
→ →
∆Es = – µ s B = – µs B cos θ
→ →
where θ is the angle between µ s and B . The
→
quantity µs cos θ is the component of µ s parallel
to B. In the case of spin magnetic moment of the
electron this component is
µBz = µB , the Bohr magneton.
µs cos θ = µsz = µB
= ±
()
e
h
2π
2m
∴ ∆Es = ± µB B
Depending upon the orientation of its spin
vector, the energy of an electron in a given atomic
quantum state will be higher or lower by µ B B
than its energy in the absence of spin-orbit
coupling. The result is the splitting of every
quantum state two separate substates. Hence every
spectral line splits into two components (lines.)
Notice that there is one exception. When an
electron is in the s state, l = 0 and as a result the
magnetic field at the site of the electron is zero.
Since B = 0, ∆Es must also be zero and there is no
fine splitting for this state.
λ = 6562·8 Å)
Fig. (a) Fine structure of Hα line (λ
λ = 4861 Å)
Fig. (b) Fine structure of Hβ line (λ
564 | Physical Sciences
λ = 4686 Å)
Fig. (c) Fine structure of helium (He) line (λ
Fine Structure of Sodium D-Line
Ten out of the eleven electrons of the normal
sodium atom are interlocked in closed shells.
They contribute nothing to the angular momentum
of the atom. We have to consider only the states
of eleventh optical electron in discussing
the spectrum of neutral sodium atom.
The D-lines belongs to principal series. Lines
of the principal series are due to the transition
from P state to the S state. For the upper P state,
3 1
L = 1, J = L ± S = or . Hence, the two
2 2
possible terms are 2P3/2 and 2P1/2 for the lower S
1 nuclear spin because of which these will be an
state, L = 0, J = , so that only one term 2S 1/2 is
2
possible. The figure given below shows the two
transitions possible between the two terms of P
state and the single term of S state. These are
(i) 2P1/2 → 2S1/2, which results in the D1 line
of sodium having wavelength 5896 Å and
(ii) 2P3/2 → 2S1/2, which results in the D2 line
of sodium having wavelength 5890 Å.
L J State
1 3/2 2P3/2
1 1/2 2P1/2
D1 D2
0 1/2 2S1/2
Now applying the selection rules ∆L = ±1 and
∆J = ±1, 0 (excluding 0 → 0), both the transitions
are allowed. This explains the doublet fine struc-
ture of sodium D line.
Hyperfine Structure
Atomic nuclei have radii of order of 10 –4 Å
which are small compared with typical distance of
an electron from the nucleus (∼ 1Å). The nuclei
are also much heavier than electrons. It is, there-
fore, a very good approximation to consider the
nuclei to be positive point charge of infinite mass.
However, high precision experiments reveal the
existence of tiny effects of the electronic energy
levels, which cannot be explained if the nuclei are
considered to be point charges of infinite mass.
These effects were first observed by Michelson in
1891 and Febry and Perot in 1897. They are
called Hyperfine effects. Because they produce a
shift in the electron’s energy levels which are
usually much smaller than those corresponding to
fine structure. The hyperfine components of spec-
tral radiations are observed using a high resolution
interferometer. The wavelength differences in the
hyperfine structure of single isotopes are very
much smaller the electron spin fine structure.
Pauli attributed the hyperfine structure to the
angular momentum of the nucleus of the atom.
The assumption of orbital and spin angular
momentum of the electron could explain com-
pletely the fine structure of spectral lines. The
hyperfine structure is caused by the properties of
the nucleus. The influence of the nucleus is
mainly due to (i) The isotope effect and (ii) The
intrinsic angular momentum.
Zeeman Effect
The splitting of spectral lines in the presence
of magnetic field is known as Zeeman effect.
Zeeman in 1896 observed that when a light source
is placed in an external magnetic field, the spectral
lines it emits are split into several polarised com-
ponents. For the fields tenth fraction of a tesla this
splitting is proportional to the strength of magnetic
field. A singlet spectral line viewed normal to the
field is split into three plane polarised components.
A central unshifted line with the electric vector
vibrating parallel to the field (called π component)
and two other lines equally spaced each on either
side with electric vector perpendicular to the field
Any singlet
NO FIELD
NORMAL TRIPLET
σ π σ
Normal Triplet
 Physical Sciences | 565

(called σ components). This combination is called Photoelectric Effect

normal triplet and the effect is called normal When a beam of radiation of frequency in the
Zeeman effect. blue or ultraviolet (UV) region falls on a metal
The fine structure components of a multiplet plate, slow moving electrons emitted from the
spectral line however, show a complex Zeeman metal surface. The electron emitted are called
pattern. For example, the D1 and D 2 components ‘photoelectrons’. The effect can also be observed
of sodium yellow doublet give four and six line from non-metals and gases, provided that radia-
respectively in Zeeman pattern. This is anomalous tion of appropriate frequency is used.
Zeeman effect. The Zeeman splitting is smaller
than the fine structure splitting. Experimental Arrangement
Sodium principal doublet Quartz window
Vacuum
NO FIELD C
Incident
A photons
MAGNETIC FIELD i
PRESENT
i
A
σπ π σ σ σ π πσ σ Variable resistance
Anomalous Pattern

Stark Effect A Anode or collector

An electric field can also, like a magnetic C Cathode
field, modify the quantum states of an atom,
separating states which are initially degenerate. Observations—(i) The emission of photo-
This effect was discovered in 1913 by Stark, who electrons cannot occur if the frequency of the
observed the splitting of the Balmer lines of incident light is below a certain value ν0 (called
hydrogen. the threshold frequency), however strong the
intensity of light, ν0 is different for different
When a hydrogen atom is placed in a
→ surface.
homogeneous electric field E directed along Z (ii) The maximum energy of photoelectrons
axis, a perturbation energy, –eEz, is added to the increases linearly as the frequency of the incident
coulomb energy in the potential energy term of light increases.
the Schrodinger equation.
It is seen that, when E is weak enough for its
effect to be smaller than the spin-orbit interaction,
each level is split into (2n – 1) components where A B C
n is principal quantum number.
Energy
For weak fields, the Stark splitting for
hydrogen atom lines leads to symmetrical pattern
and the line spacing are proportional to E. This is
called the linear or first order Stark effect.
When E exceeds 107 V/m there are shifts in
the line patterns which are proportional to E2 and
ν0 (A) ν0 (B) ν0 (C)
in this case, the effect is called quadratic or Frequency of ν
radiation
second order Stark effect.
A qualitative explanation is that under the A, B and C show energy Vs. frequency for
effect of strong electric field, the atom became three metals say A, B and C.
polarised with an induced dipole moment propor- If W is the work function for a particular
tional to E. The torque on such a dipole become metal, then
proportional to E2. W = hν0
566 | Physical Sciences

The maximum kinetic energy of electron is Although photoelectric effect provide good
equal to the energy of light minus the work support to quantum theory of light yet the complete
function, i.e., and concrete explanation can be given by Compton
effect.
[ ] 1 2
2
mv
max
= hν – hν0
Compton Effect
= h(ν – ν0) When a monochromatic beam of X-rays is
This is known as Einstein’s photoelectric scattered by a substance, the scattered radiation
equation. contains original radiation’s wavelength as well as
(iii) The maximum energy of the photo- longer wavelength also. This effect was explained
electrons is independent of the intensity of on the basis of quantum theory of radiation.
radiation. An X-ray photon striking an electron
(iv) The number of photoelectrons emitted (assumed to be initially at rest in the laboratory
per second (photoelectric current) varies directly coordinate system) and being scattered away from
as the intensity of incident light. its original direction of motion while the electron
receives an impulse and begins to move. In the
collision the photon may be regarded as having
lost an amount of energy that is same as the
kinetic energy K gained by the electron, though
actually separate photons may be involved. If the
Current
initial photon has the frequency ν and the
scattered photon has the lower frequency ν ′,
where
hν – hν′ = K
Vο
Retarding potential
If V0 is the retarding potential, then
1
eV 0 = [mv2]max
2
or eV 0 = h(ν – ν0).
(v) Photoelectrons are emitted almost instan- By using laws of conservation of momentum
taneously, however poor the intensity of the
and energy (In two directions) we can calculate
incident light. The photoelectric current is inde- the shift in wavelength. [Remember, in Photo-
pendent of the frequency of the light used. electric effect we consider the operation with
energy only while in Compton effect energy and
momentum both are considered].
We get the shift,
h
∆λ = λ′ – λ = (1 – cos φ)
m0 c
where
λ → incident wavelength,
λ′ → scattered wavelength,
m0 → rest mass of electron,
and φ → angle at which photon is scattered.
In the above diagram : If we define λc as Compton wavelength
π
ν1 > ν2 > ν3 and light intensity = constant
the photoelectric effect can be explained by
λc =
h
m0 c ( )φ=
2
quantum theory of light only. λc = 2·426 × 10–12 m

Then the maximum possible wavelength
change corresponding to φ = 180° is
λ′ – λ = 2·426 × 10–12 {1 – (–1)}
= 4·852 × 10–12 m
Change of this magnitude of this order or less
are readily observable only in X-rays.
This explains why Compton effect is not
observed in visible light as the shift is less than
0·01 per cent of the initial wavelength.
If electron is not loosely bound to parent
atom, the entire atom is recoiled and the shift is
very small because m 0 increases on account of
total mass of atom.
A free electron can not absorb all the energy of a
photon and still conserve both energy and momen-
tum.
We can, hereby, also conclude that the photo-
electric effect cannot occur for completely free
electrons. Whereas in Compton effect, the electron
can be free because it absorbs only a part of energy
of incident photon.
X-ray Spectra
The X-ray spectra is the result of reverse
photoelectric effect with the electron’s kinetic
energy being transformed into photon energy h ν.
The optical spectra and X-ray spectra differs from
each other in terms of position of ejected electrons.
The transition of electrons from outer orbits
involves small energy hence the resulting photon
lies in or near the visible region, whereas in case
of X-ray spectra the tightly bound inner electrons
are ejected out by highly energised striking elec-
trons resulting X-ray series. The energy level
diagram for a heavy atom is shown below :
Mγ
Lβ
Lα
Lγ
Lδ
β
αM
M
K α
K
β
Kγ
Kδ
Kε
K L M N OP
The continuous spectrum graph drawn
between relative intensity (Energy density) of the
X-rays and the frequency show sharp maxima
lines superimposed on the continuous spectrum.
These sharp lines were due to the impinging
Physical Sciences | 567
electrons penetrating the atoms of the target and
dislodging the electrons in the K, L, M etc; inner
orbits of the target atoms. When electrons in the
outer shells at the higher level jump into the
vacancies caused in the inner shells, an X-ray
quantum is emitted which gives the characteristic
maxima lines in the continuous spectrum.
The X-ray spectra arises due to vacancies
caused in the inner shells K, L, M … etc; of heavy
atoms and consequently emission of X-ray quanta
by filling of the vacancies by electrons is the outer
shell.
Moseley’s Law
Moseley’s discovered that in the X-ray spectra
the square root of frequency of the Kα line was
proportional to the atomic number of the emitter
target. In general he deduced,
ν = A (Z – σ)2
here A and σ are constants and Z is the atomic
number of target nucleus.
Moseley’s law can be used to establish the
sequence of elements in the periodic table also in the
prediction of yet to discovered elements.
The continuous X-ray spectrum has a short
wavelength limit (see figure) and the cut-off
wavelength can be determined by
eV = hνmax
hc
⇒ eV =
λmin
This formula is known as Duane-Hunt
formula, it is very useful in the determination of
Planck’s constant h.
λmin is
independent
Relative Intensity
of the target
material
50 KV
40 KV
30 KV
20 KV
Wavelength (λ) Å
[Continuous X-ray spectra showing λmin].
The continuous X-ray spectrum is generated
when electrons of relatively low energy fall on the
target atom. As the energy of striking electrons
568 | Physical Sciences
increases the characteristic lines starts predomina-
ting over continuous spectrum.
Screening Effect
For an electron which in orbiting in L shell,
the K shell electron behave as if it is merged into
the nucleus thereby reducing the atomic number
by two. The reduction procedure of nuclear charge
by electron is known as screening effect. In the
same fashion for a M electron K and L shell
electrons screen the nuclear charge. Therefore, the
energy term value of the system is written as
R(Z – b)2
T =
n2
here (Z – b) denotes the effective atomic number
of element.
X-ray Absorption Edges
When a beam of X-rays passes through a
material, some of the photons may interact with
the atoms of the material, causing the photons to
be removed from the beam. The primary interac-
tion processes responsible for the reduction of
the intensity of any photon are the photoelectric
effect, compton scattering and pair production.
Because X-ray have energies in the 1–100 keV
range they cannot produce electron-positron pairs,
which require energies in excess of 1000 keV.
Therefore, the intensity of X-ray beam will be
reduced by only the first two of the above
processes, with the photoelectric effect being the
dominant mechanism.
The intensity I of a monochromatic X-ray
beam after it has penetrated a distance x in a target
material is given by
I = I0 exp (– µx )
where I0 is the intensity of the incident beam and
µ is the absorption coefficient of the material. The
quantity µ depends on both the target atoms and
the energy of X-ray photons.
Suppose for a given target material, that µ is
measured as a function of the incident X-ray
energy. As this energy increases, the absorption
coefficient decreases because the higher energy
photons are less likely to produce photoelectron or
undergo Compton scattering. This decrease
continues until the X-ray energy just equals the
binding energy of one of the core-electrons. At
this point, more electrons suddenly become avail-
able for photoelectric emission, causing a marked
decrease is the transmitted X-ray intensity or
equivalently, a sudden increase in the value of the
absorption coefficient. This sharp increase in µ
happens at the binding energies of each of the
core electrons, resulting in absorption edges
shown in the figure given below (a). Measurement
of the energies of the K, L, … absorption edges
thus serves to determine the binding energies of
the corresponding core electrons.
With the exception of the K edge, each
absorption edge actually consists of a number of
closely spaced peaks corresponding to the fine
structure of the energy levels figure (b).
LI
Absorption coefficient µ
LII
LIII
Photon energy hν
(b)
Characteristics of Molecular Spectra
The molecular spectra can be divided into
three categories :
(a) Rotation spectra—A rigid molecules
containing electric charges such that it possessed
an electric moment. If it were to rotate, it would
emit radiation. Conversely if radiation falling
upon such a molecule sets it into rotation, energy
of the radiation would be absorbed.
Spectral lines corresponding to the rotational
spectra are observed in FIR region.
(b) Vibration rotation spectra—Non-rigid
molecules contains atom which vibrate under
elastic forces about equilibrium positions, radia-
tions would be emitted. The radiation frequency is
 Physical Sciences | 569

same of atomic vibration frequency. Also the 3 → 12B

rotation of molecules at the same time gives rise 2 → 6B
to fine structure of lines, i.e., the lines have 1 → 2B
frequencies lesser of greater than the atomic J = 0 → 0
vibration frequency. The spectral lines of this
The difference between the energy levels can
series lies in IR region.
be calculated :
(c) Electronic spectra—The electrons in the F(J) for J = 0 → 1
molecules vibrate and hence radiations emits. The = 2B cm–1
vibration and rotation of atom affects the F(J) for J = 1 → 2
electronic radiations. The lines of this spectra falls
= 4B cm–1
in visible and UV region.
F(J) for J = 2 → 3
Rotational Spectra = 6B cm–1
A di-atomic molecule having an internuclear and so on.
distance r0 is shown in the figure. There are only In general to raise the molecule from state J
two mass points and the line joining the two is an to J + 1,
axis of symmetry. The molecule rotates end over F(J)J → J + 1 = B(J + 1) (J + 2) – B(J) (J + 1)
end about a point c, which is the centre of mass. = 2B (J + 1) cm–1
r0 (for J = 0, 1, 2, …)
Thus, the stepwise raising of the rotational
C energy levels results in an absorption spectrum
m1 m2 consisting of lines at 2B, 4B, 6B, … 2B(J + 1)
cm –1 while a similar lowering would result in
identical emission spectrum.
r1 r2
Frank Condon Principle
m1 r 1 = m2 r 2
The different types of intensity distributions
The moment of inertia about c is given by of vibrational electronic spectra can be explained
I = m1 r 1 2 + m2 r 2 2 on the basis of Frank Condon principle whose
⇒ I = µr0 2 basic semi-classical idea is as follows :
“An electronic transition in a molecules takes
where µ is the reduced mass place so rapidly compared to the vibrational
m1 m2 motion of the nuclei that the instantaneous inter-
µ = and r0 = r1 + r2 .
m1 + m2 nuclear distance and the velocity of the nuclei can
The allowed energy levels can be given by be considered remaining unchanged during the
solution of Schrodinger’s equation for a rigid electronic transition.”
rotator Raman Spectra
h2
EJ = J (J + 1), J = 0, 1, 2, … When a monochromatic radiation is scattered
2I by a transparent liquid then the scattered light not
here J is called the rotational quantum number. only consist of the radiation of incident frequency
This energy EJ can be converted to the units of but also the radiations of frequencies above ε-
term values, below that of incident beam frequency and it is
E called Raman effect.
h
F(J) = J = 2 J (J + 1) If νi is the frequency of incident radiation and
hc 8π Ic
νs is frequency of light scattered by a molecular
= B J (J + 1)
species then Raman shift be
h
here B = 2 is rotational constant for molecule. ∆ν = νi – νs
8π Ic This difference is characteristic of the subs-
These energy levels can be shown as follows : tance producing scattering
5 → 30B ∆ν is positive → Condition of Stoke’s line
4 → 20B ∆ν is negative → Anti Stoke’s line
570 | Physical Sciences
Laser
An excited atom return to its ground state by
the emission of radiation. This energy emission by
the atom can be (i) Spontaneous or (ii) Stimulated.
Spontaneous emission can occur in high
frequency transition (IR visible and UV). How-
ever stimulated or induced emission is a reso-
nance phenomenon. An excited state is droped to
its ground state only when a photon (or radiation)
of the same frequency ν interacts with the system.
This is called the stimulated emission which is
responsible for the functioning of laser. These
radiations are highly coherent, highly intense and
precisely parallel in nature.
The straight forward mechanism is given
below :
Here E0 → Ground state
E1 → Metastable state
E2 → Excited state.
and various step represents :
Step I. Optical pumping
Step II. Rapid transition to metastable state
by spontaneous transition.
Step III. Population inversion.
Step IV. Induced emission and secondary
photon.
Social Inclusion Principle
When there are n identical boson present in a
given quantum state, the probability that one more
particle will enter the same state is enhanced by a
factor (n + 1) bigger the gathering, better the popula-
tion. This highly social nature of boson is known as
social inclusion principle.
The stimulated emission is a general conse-
quences of this principle.
Electron Spin Resonance (ESR)
Electron spin resonance is a branch of
absorption spectroscopy in which radiation having
frequency in the microwave region is absorbed by
paramagnetic substances to induce transitions
between magnetic energy levels of electrons with
unpaired spins.
Some atoms and molecules (e.g., N O 2 , O 2
and NO) contains one or more electrons with
unpaired spins. These substances show E.S.R.
spectra. The intensity, width, position and multi-
plet structure of spectral lines are of importance in
E.S.R. spectroscopy. The intensity of an E.S.R.
absorption is proportional to the concentration of
the free radicals or paramagnetic material present.
The width of an E.S.R. resonance depends on the
relaxation time of the spin state
under study. Of the two possible
relaxation processes, the spin-spin
interaction is usually very efficient
and gives a relaxation time of
10 –6 – 10–8 sec the spin-lattice
relaxation is efficient at room
temperature (some 10 –6 sec). For
most samples, we could choose
10–7 sec as a typical relaxation
time. Using Heis-enberg uncertainty
relation, frequency uncertainty
(line width) is found to be = 1
MHz.
Theory of ESR
The presence of an unpaired electron in a
molecule or ion allowes energy levels to be
produced from the interaction of the magnetic
moment of the electron with an applied magnetic
1
field. For an electron of spin S = , the spin
2
angular momentum quantum number can have
1
values of m s = ± . In the absence of magnetic
2
1 1
field, the two values of m s i.e., ± and – will
2 2
give rise to a double degenerate spin energy
state. If a magnetic field is applied, this
degeneracy is removed and this leads applied field
by amount + µeH and – µeH. Here µe is the
magnetic moment of the spinning electron and H
the magnetic field acting on the unpaired electron
(see figure).
 Physical Sciences | 571

microwave region. When the absorption occur,

External field applied
In ESR, a transition between the two different
electron spin energy states occur upon absorption
of a quantum of radiation of frequency ν in
the following relation holds good.
2µeH = hν
This equation holds good for a free electron.
The energy (∆E) of transition is more accu-
rately given by
∆E = hν = gβH
here h is Planck’s constant β is Bohr magneton
and g is Lande splitting factor.

Solved Questions
Q. 1. If ν H and ν D are the frequencies of From these equations
the transitions in hydrogen and deuterium
νH MH  MD + m 
MD  MH + m 
atom respectively corresponding to same set of =
quantum numbers, then show that the mass of νD
electron, m is given by
MD + m M D νH
MHMD (ν νD – νH) or =
m = MH + m M H νD
MDνH – MHνD
here MH and MD are the masses of hydrogen Subtracting (1) from both the sides
and deuterium nuclei respectively. M D – MH MDνH – MHνD
=
Solution : Deuterium is the isotope of MH + m M H νD
hydrogen hence, we can use Z = 1 for both the
elements, from Bohr’s theory MHνD (MD – MH)
or MH + m =
MDνH – MHνD
νH = R H ( )
1 1
–
nf 2 ni 2
…(1)
MHνD (MD – MH)
or m = – MH
MDνH – MHνD
and νD = R D (n1 – n1 )
f
2
i
2 …(2)
MHMD (νD – νH)
⇒ m =
Also RH =
R∞ MDνH – MHνD
m
1+ Q. 2. Show that all eigen functions which
mH
are solutions of time independent Schrodinger’s
R∞ equation having spherically symmetrical poten-
and RD =
m tial have definite parities, either even or odd.
1+
mD
Solution : A spherically symmetric potential
m is one which has the same value at points (x, y, z)
1+
RH MD and (–x, –y, –z).
or =
RD m
1+ The time independent Schrodinger equation is
MH
8π2m
RH MH  MD + m  ∇2 ψ + h2 (E – V)ψ = 0
MD  MH + m 
∴ = …(3)
RD
Another form of it will be
Dividing equation (1) by equation (2)
RH ν
 h2 
= H  – 8π2m ∇2 + V ψ(x, y, z) = Eψ(x, y, z)
RD νD

CONDENSED MATTER PHYSICS
Crystallography
The science of crystallography is concerned
with the classification of all types of crystal
structures and the determination of the actual
structure of crystalline solids.
Amorphous and Crystalline Solids
The intermolecular forces bind a large number
of atoms into a form which has a volume and
definite shape of its own. Such a system of atoms
which preserve their volumes and shapes unless
subject to large external force are called solids.
Solids may be classified into two types :
(1) Amorphous solids and (2) Crystalline
solids.
Amorphous solids are termed as the super
cooled liquids having small structural units with a
short ranged order. These solids are those which
lack the regular arrangement of atoms or molecules
and hence a short range order in their structure.
Amorphous solids have a stiffness and rigidity due
to their highly enhanced viscosity and the small
structural units of such solids remain fixed in their
positions while they are free to roam all the
volume of the liquid. Another important property
of these solids is that the bonds between atoms
among them are not equally strong. Hence, on
rising the temperature, the weakest bonds are
ruptured first at lower temperature and the
stronger ones at higher temperatures. Due to this
fact an amorphous solids has no sharp melting
point because it starts softening gradually before
finally melting. Its examples are pitch, glass and a
few plastic substance.
All materials are composed of atoms or
molecules. In crystals the atomic array is periodic
and hence a long range order in their structures. In
crystalline solids all the bonds are equally strong
due to long range symmetry. Hence, when heated
Unit-9
all of them get ruptured at a unique temperature.
That’s the reason, therefore, known by sharp
melting points. The examples of crystalline solids
are sugar, quartz, rock salt, calcite etc. A crystal
may also be defined as solid in which the atoms
are regularly arranged in such a way that the
distance between the successive atomic layers is
of the order of few Å.
Elements of Symmetry
Let us suppose an array of points in such a
way that the surrounding about one point is
identical with the surrounding about any other
point. Such an array of point in two dimensions as
shown in given figure is called the plane lattice.
To construct the two dimensional lattice, let we
choose any two convenient axes OA and OB and
then the points lie at equal intervals ‘a’ along OA
and at equal intervals ‘b’ along OB.
B
b
A
O a
If we extend this array of points in three
dimensions then the array of points is called
space-lattice. In order to construct the space-lattice
or 3d-lattice, the points are arranged at equal
intervals c in the third direction also. The position
vector of any lattice in two dimensional lattice can
be expressed by choosing any other lattice as
origin. Hence,
→ → →
T = n1 a + n2 b
628 | Physical Sciences

Here ‘a’ and ‘b’ are called the primitives and
n1 and n2 are integral values that show the number
of lattice points along OA and OB axes. T is
translational operator.
Similarly, the arrangement of points in three
dimensions can be shown by the translational
operator T as
→ → → →
T = n 1 a + n 2 b + n3 c
where the primitives a, b and c are taken along the
crystal axes OA, OB and OC.
The parallelopiped formed by using the trans-
lations a, b, c as concurrent edges, is called the
unit cell of the space-lattice. The angles between
primitives (b, c), (c, a) and (a, b) are shown by α,
β and γ-respectively.
Therefore, a lattice is defined as a parallel net
like arrangement of points provided the environ-
ment about any point is identical with surrounding
about any other point.
crystals and the corresponding Bravais lattices
with their properties are given below—
1. Cubic crystals—In these crystals, axes are
mutually perpendicular that means (α = β = γ =
90°) and the repetitive interval is the same along
the three axes (a = b = c). In cubic crystal, the
lattice may be simple body-centred or face
centred. In cubic lattice, the lattice points lie only
at the corners of the unit cell, while in the body
centred lattice, the lattice points are situated at the
corners and also at the intersection on the body
diagonals of the unit cell. In face-centred lattice
the lattice points are situated at the corners as well
as at the centres of all the six faces of the unit cell.
NaCl crystal has face centred cubic lattice while
CsCl have simple cubic lattice.

2. Tetrogonal crystals—These crystals may

be simple or body centred. They have axes
perpendicular to one another (α = β = γ = 90°) and
have repetitive intervals along two axes same but
interval along the third axis is different, it means
Simple Cubic Lattice that (a = b ≠ c).
Bravais discovered the fact that there could
only 14 different types of lattices under the seven
systems of the crystals. These lattices are called
Bravais lattices. For a cubic system three types of
lattices are possible which are—
(a) There is one lattice point at each of the
corners of the unit cell. This type of cell, is called
simple or primitive cubic cell (P) of the system.
3. Orthorhombic crystals—These crystals
(b) There is one lattice point at each of the may be simple, base centred, body centred or
eight corners and one lattice point at the centre of face centered. These crystals have perpendicular
the cubic cell. This type of cell is called body axis to one another (α = β = γ = 90°) but the
centred cubic cell (I). repetitive intervals are different along different
(c) There is one lattice point at each of the axes (a ≠ b ≠ c).
eight corners and one lattice point at the centre of
each of the six face centred cubic cell (F).
Seven Crystal Systems
There are seven types of crystals depending
upon their axial ratios (a : b : c) and angles
between them (α, β, γ). Bravais told that they can
give rise to 14 types of space lattices. These
 Physical Sciences | 629

4. Trigonal crystals—The trigonal lattice is

only simple. The angles between each pair of
crystals axes are the same but different from 90°
(α = β = γ ≠ 90°).

5. Monoclinic crystals—Monoclinic lattice

may be simple and base-centred. The repetitive
intervals are different along all three axes (a ≠ b
≠ c) and two of the crystals axes are not. Sodium Chloride Structure
NaCl compound has the unit cell sketched in
given figure.

6. Hexagonal crystals—In these crystals

(α = β = 90°, γ = 120°) and (a = b ≠ c)

Fig. : NaCl crystals structure

It consists of two face centred cubic sub-
lattices, one of Na ions having its origin at the
point O, O, O, the other of Cl ions having its origin
a
midway along a cube edge say; at point , O, O.
2
This lattice is simple cubic if the difference bet-
ween the Na and Cl ion positions is ignored.

Cesium Chloride Structure

7. Triclinic crystals—In these crystals (α ≠ β
≠ γ) and (a ≠ b ≠ c).
The lattice points of this compound are two
interpenetrating simple cubic lattices, the corner of
one sub-lattice is the body centre of the other. One
of the sub-lattice is occupied by Cs ions, the other
by Cl ions. The resultant structure, is sketched in
figure given below :

Crystals Structure
The unit assembly of atoms is called the
basis. The crystals structure may be considered to
be formed by the addition of a unit assembly or
basis of atoms to every lattice point of the space-
lattice. The lattice and crystals are distinguished in
given figure. Hence, it is clear that
Crystals structure = Lattice + Basis
630 | Physical Sciences
Zinc Blende Structure
A new structure occurs if the two face-
centred cubic sub-lattice are occupied by different
elements. The resultant structure termed the zinc
blende structure is shown in given figure. Impor-
tant compounds which have this structure are the
semiconductor in Sb and GaAs also ZnS and CuCl
and many other compounds.
Zinc blende structure
Determination of Crystal Structure
The general symmetry characteristics of
crystals can be inferred from a combination of
measurements of both the visible, external faces of
crystals and bulk physical properties. Symmetries
determined in this way do not however, fix the
positions of atoms on the crystal lattice nor do
they permit measurement of distance between
atoms. The atoms in a crystal are simply too close
together to be resolved by such observations.
Visible light with a wavelength of about 5000A
does not allow the resolution of the positions of
the atoms themselves. Since, inter-atomic distan-
ces are only a few angstrom units, the appropriate
radiation is the X-ray.
Bragg’s Law
The conditions for diffraction in actual
crystals have been conveniently systematized by
Bragg who reduced the nomenclature of diffrac-
tion to one of specular reflection off the lattice
planes. Any three dimensional lattice has an
infinite number of sets of equally spaced atom
planes. Bragg has shown that any diffracted ray
can be regarded as if it were reflected from one of
these systems o f planes much a s through it were
reflected from a mirror parallel to planes.
However, the given wave can be reflected only if
it satisfies the equation
nλ = d sin θ
Here n = An integer
λ = Wavelength of the radiation
θ = Angle of the incident radiation
d = Distance between adjacent planes
The geometry of this arrangement is shown in
figure :
The above relation shows that why radiations
in X-ray region of the spectrum are most useful
for crystal analysis. The interatomic spacing of
solids is about 2Å. Since, sin θ cannot be greater
than 1, first order Bragg reflections off neighbour-
ing planes of a set, require λ should be 2Å or less.
Therefore, X-rays with wavelengths less than 2Å
are most useful in studying crystals.
Debye and Scherrer Method or Powder
Method
This method is used in crystal structure
analysis to determine the various d-values for the
crystals of a given material.
In the powder, the millions of micro-crystal
have all possible orientation in space. When the
collimated X-ray beam falls upon the powder, for
each set of d and n, there is a particular value of θ,
which satisfies Bragg’s equation. For this value of
θ, the beam appears at the corresponding 2θ
deviation. For various sets of d and n, we obtain
various cones of ray’s. One such cone with a
semi-vertical angle 2θ is shown in given fig. (a)
and the pattern recorded on the photographic film
is shown in fig. (b).
 Physical Sciences | 631

If l is the distance between the arcs of the drift velocity of free electrons is superimposed
same circle on the photographic film and R is the over their random velocity. Hence, the random
radius of the camera chamber. Then we get, motion of free electrons is modified by the
l application of external random speed of free
2θrad = electrons.
R
l 7. There is a loss of kinetic energy in the
or θrad = collision of free electrons with the positive ion
2R
fixed in the presence of external eletric field. Such
collisions are called inelastic collision. Just after
the inelastic collision, the electron momentarily
will have random motion. Hence, in the presence
of electric field, the process of inelastic collisions
of free electrons with the positive ions of the
lattice continues.
8. The average distance traversed by a free
(a)
electron between two successive collisions with
the positive ion is known as mean free path
denoted by λ.
(b) ne2τ
∴ Electrical conductivity σ =
Hence, by measuring the diameter of these 2m
arcs (2l) and using the radius R of the camera Here n is the total number of free electrons in
chamber, we can measure θ. the metal, e is charge of electron, τ is average bet-
‘‘On the basis of free electron theory ween collision and m is the mass of free electron.
electrical conductivity of metals’’—In 1900, If λ is the mean free path and vF the speed of
Drude was first man to put forward a simple free electrons, then average time τ between
theoretical explanation of electrical conductivity collision is given by
of metals. Later on Lorentz modified this theory.
The fundamental assumptions of Lorentz-Drude λ
τ =
theory are— vF
1. There are a large number of free electrons neλ
So, σ =
in a metal and these electrons are free to move 2mv F
above the whole volume of the metal. The drift velocity per unit electric field is
2. The assembly of free electrons in a metal is called mobility. So,
called electron gas. The electrical or thermal V
conductivity of metals is solely due to these free µ = d
E
electrons.
eEτ eτ
3. The average kinetic energy of an electron is = /E=
2m 2m
()3
2
kT where k is Boltzmann’s constant and T is Therefore, electrical conductivity
ne2τ
absolute temperature. σ =
2m
4. The current density vector is zero in the
= neµ
absence of electric field.
5. The collisions among the electrons From above relation it is clear that at a given
temperature the only factor which varies from one
themselves have no practical contribution to the
conductivity of metals. to other is ‘n’ (the number of free electrons per
unit volume). Hence, the great difference in the
6. When an external electric field is applied, conductivities of different metals are due to the
the electrons drift slowly with some average difference in the number of free electrons per unit
velocity; it is called average drift velocity. The volume.
632 | Physical Sciences
Heat Capacity of Electron Gas
The specific heat of metal has the
contribution from the atoms as well as the free
electrons. So,
C V (metals) = C V (atomic) + CV (electronic)
The atomic contribution to specific heat
referred to 1 gm. atom of metal is
C V (atomic) = 3R
(R = gas constant for 1gm atom)
3
and C V (electronic) = R
2
Hence, classically the specific heat referred to
1 gm atom of metal is
3 9
C V = 3R + R = R
2 2
It is too high as compared to the experimental
value 3R. This specific heat anomaly was solved
by Fermi-Dirac statistics.
According to him the electronic contribution
to the specific heat per gm atom of metals at low
temperature is
CV (electronic) = λT
1 . nk2π2
Here λ =
2 εF (0)
and n = Number of electron/m3
Obviously, the electronic contribution to the
specific heat of a metal at low temperatures is
proportional to the temperature and vanishes at
absolute zero. This expression agrees with the
experimental observations.
Again from the Debye’s theory, the specific
heat at extremely low temperature is proportional
to T 3 . Therefore, the specific heat of metals at low
temperature can be expressed as—
C V = βT 3 + λT
where the first term refers to the atomic con-
tribution to the specific heat according to Debye’s
theory and the second term refers to the electronic
contribution to the specific heat according to
Fermi-Dirac statistics. At very low temperatures
the atomic specific heat becomes small and the
electronic specific heat becomes relatively high.
Thus, the Fermi-Dirac statistics has solved the
problem of specific heat of metals.
Electron Motion in a Periodic Potential
The simplest quantum mechanical view of the
potential experienced by an electron in passing
through the crystal is one which is perfectly
periodic as shown in given fig. The distribution of
potential between the nuclei is a series of humps
and increases at the boundary of the solid. The
periodicity of the potential is the result of coulomb
interaction between the moving electron and the
periodic charge distribution arising from the posi-
tive ions situated on the lattice sites. Thus, the
periodicity of the potential is equal to the period-
icity of the lattice.
Distance (x) through crystal
The solution of Schrodinger wave equation
for the single electron in this one dimensional
periodic crystal potential provides the set of one
electron states which the electron may occupy.
The problem was simplified by replacing the
periodic potential with rectangular barries as given
in figure. It is called Kroning Penny model.
The Schrodinger wave equation to be solved
is
d2 ψ (x) 2m
+ h 2 [E – V (x)] ψ (x) = 0 …(1)
dx2
where V(x) = 0 for regions of type I and V(x) = V0
for regions of type II. The effect of the periodic
potential on the wave function is to change the
wave function of the electron eikx to
ψ (x) = u (x) e ikx,

√ 2m (E – V (x))
K = …(2)
h
The modulating amplitude u (x) is almost
unity at the centre of the well. Hence, u (x) repeats
from lattice to lattice and satisfies the condition
u (x) = u (x + a) …(3)
Here a = b + c is the lattice spacing.
 Physical Sciences | 633

Equations (2) and (3) together are known as solid material, then there will be 50 levels of
Bloch’s theorem which is a mathematical state- slightly different energies, i.e., 50 energy level
ment for the form of one electron wave function diagrams would be superimposed on each other.
for a perfectly periodic potential. The condition Consequently, the split energy levels are almost
for the solution of wave equation to exist in the continuous and are said to form an energy band.
form is : The energy bands in a solid correspond to the
ρ sin αa energy levels in an atom. An electron in a solid
cos Ka = + cos αa …(4)
αa can have only those discrete energies that lie
Here ρ, α and β are constants given by within these energy bands. These bands are, there-
fore, called allowed energy bands. These energy

√ 2mE ,
α = h
bands are, in general, separated by some gaps
which have no allowed energy levels. These gaps

β =
√
 2m (V0 – E)
…(5)
are known as forbidden energy bands. Energy
h band occupied by the valence electrons is called
valence band and is obviously the highest occupied
mV0ba band. Electrons which have left the valence band
and ρ = h2 …(6)
are called conduction electrons and are only
weakly bound to the nucleus. The bands occupied
Energy Bands in Solids by these electrons is called the conduction bands.
For a single isolated atom, there are discrete Thus, the band beyond forbidden band is called
energy levels 1s, 2s, 3s, …… that can be occupied conduction band, into which, when the electrons
by the electrons of atom. All the atoms of a solid, pass, they can move freely.
if assumed isolated from one another, can have
completely identical electronic schemes of their Distinction between Conductors,
energy levels. Then the electrons fill the levels in
each atom independently. Insulators and Semiconductors
An isolated atom are brought together to form Every solid has its own characteristic energy
a solid, various interactions occur between band structure. This variation in band structure is
neighbouring atoms. As a result of this interaction, responsible for the wide range of electrical charac-
the higher energy levels are considerably affected, teristics observed in various materials.
i.e., the energy levels of the outer shells are The electrical properties of the materials can
slightly altered without violating Pauli’s exclusion be explained by the concept of energy bands.
principle. There will be splitting of a single energy Depending on the nature of band occupation by
levels of an isolated atom into a large number electrons and on the width of forbidden bands, all
of energy levels as shown in fig. (a). solids can be classified as conductors, insulators
Valence Band
and semiconductors, which are given below :
Forbidden Band
1. Conductors (metals)—The arrangement
of energy bands in a conductor (metal) is shown in
Conduction Band
fig. (a). The characteristics of the formation of this
Split levels band is that in conductors, the valence band is
3s
2p
of outer either partially filled or the next allowed empty
shells band (conduction band, into which electrons can
2s
2p pass) overlaps with the filled valence band; if
1s there are unoccupied states for electrons in the
2s uppermost band and these electrons are available
(a) For isolated atom (b) For atoms in a solid to carry current.
Since in a solid, many atoms (N = 10 23 Let us consider the example of sodium (Z = 11)
atom/cm3 ) are brought together, the separation which has an electronic configuration 1s2 2s2 2p6
between N sub-levels (into which each discrete 3s1 . It shows that the 1s, 2s and 2p bands in
energy level splits due to interaction) is very small. sodium are completely filled while the outermost
For example, if there are 50 atoms in a piece of (3s) band has only one electron whereas it can
634 | Physical Sciences
accommodate one more and, therefore, valence
band is half-filled. When an electric field is
applied across a piece of sodium crystal, electrons
easily acquire an additional energy, in the form of
kinetic energy to move to conduction band where
they can move freely, constituting an electric
current. Sodium and all other metals are,
therefore, good conductors of electricity.
Thus in conductors, we have overlapping
valence and conduction bands. Hence, the valence
band itself is also the conduction band.
As expected from the metallic band structures,
metals have a high electrical conductivity.
Conduction
Band
Conduction
Band
Conduction
Band
6 ev Forbidden
Overlap Band
Valence
Band
Valence Valence
Band Band
(a) Conduction (b) Insulator (c) Semi-
conductor
2. Insulators—The energy band diagram of
insulators is shown in fig. (b). A solid behaves as
an insulator if it satisfies two conditions—(a) it
has an even number of valence electrons per atom
and (b) the valence and conduction band are sepa-
rated by an energy gap large compared with kT.
Thus, diamond with four valence electrons per
atom and having a large gap of 6 eV between
valence and conduction band, is an insulator. It
has completely filled valence band and a large
(≈ 6 eV) forbidden band, i.e., electrons require
6 eV energy to jump to conduction band from
valence band. An electric field cannot give this
amount of energy to an electron in the solid.
Therefore, this type of substances behave an
insulators because no free electrons are available
to conduct the electricity. Thus in insulators, there
is—
(a) Full valence band,
(b) An empty conduction band,
(c) A large energy gap (of several eV) bet-
ween them.
3. Semiconductors—Semiconductors have a
forbidden band (energy gap) between filled
valence band and unoccupied conduction band.
Which is small as compared to the gap in case of
insulators [fig. (c)]. For example, the width of
forbidden band, which is ≈ 6 eV in diamond
(insulator), is 1·2 eV in silicon and 0·7 eV in
germanium (semiconductor). The narrow energy
gaps in Si and Ge imply that some electrons from
the completely filled lower valence band can be
thermally excited to the empty conduction band
where they are able to conduct small current on
the application of electric field to the crystal.
In a semiconductor at room temperature, there
is—
(a) A partially filled conduction band.
(b) A partially filled valence band.
(c) A very narrow energy gap (≈ 1 eV) bet-
ween them.
Free electrons
Intrinsic Semiconductors
A semiconductor in its natural pure form is
1 eV called an intrinsic semiconductor. The conduc-
tivity of an intrinsic semiconductor is solely
determined by the thermally generated carriers.
Holes
In Intrinsic Semiconductors Current
Flow
In order to understand the electrical conduc-
tion in intrinsic semiconductors, let us consider
the case of germanium. Atomic number of
germanium in 32, so it has 32 electrons—2 in the
first orbit, 8 in the second orbit, 18 in the third
orbit and remaining 4 in the outermost orbit. Thus,
germanium atom has four valence electrons, i.e., it
is a tetravalent element. Each of the four valence
electrons in a germanium atom is shared by the
valence electrons of four adjacent germanium
atoms and make four electron pairs as shown in
given figure. This type of electron pairs are known
as covalent bonds which provide the binding force
between neighbouring atoms. In this way, the
atoms behave as if their outermost orbits were
complete with eight electrons. Thus, no free
electrons are available to conduct a current
through a germanium crystal. Consequently, pure
germa-nium is actually an insulator at low
temperature (0 K). However, with increase in
temperature, a few covalent bonds are broken due
to thermal agitation. Thus, some of the electrons
become free to con-duct current. These electrons
move freely through the crystal lattice in a random
manner. On applying a suitable potential diffe-
rence across the germanium crystal, this random
motion of electrons experiences a steady drift
towards the positive electrode and a weak electric
current is produced.

When a covalent bond is broken due to
thermal energy, the removal of one electron leaves
behind an empty space or electron deficiency
known as hole. The hole acts as positively charged
particle, having a charge equal to that of an elec-
tron but of opposite sign. Since a hole is a strong
centre of attraction for the electron, the hole
attracts an electron from the neighbouring atom to
fill it. Thus, the hole is now shifted to another
place from where the electron has migrated. The
newly created, hole is filled up by another neigh-
bouring electron and so on. Thus, a hole once
created moves about in the crystal in a random
manner, just like the free electron. On applying an
electric field, the holes move in a direction
opposite to that of the valence electrons. It
constitutes a hole current.
It must be seen that in an intrinsic semicon-
ductor, the number of electrons and hole is very
small and hence the current produced in a semi-
conductor is not adequate for any useful work.
Extrinsic Semiconductors
It has been observed that if a small amount
of impurity is added to a semiconductor (generally
1 impurity atom for 108 atoms of semiconductor),
it significantly increases the conducting proper-
ties. The process of adding the impurities to a
semiconductor is known as doping and the impu-
rity that is added is called the dopant. A semi-
condcutor containing impurity atoms is called the
doped impurity on extrinsic semiconductor.
Depending upon the type of impurity added in
pure Ge on Si, the extrinsic semiconductors are
further sub-divided into two groups.
Physical Sciences | 635
(i) n-type semicondcutor
(ii) p-type semicondutor.
(i) n-type semiconductor—If a small amount
of pentavalent impurity (having five valence
electrons) is added to a pure semicondcutors, it is
known as n-type semiconductor.
In extrinsic semiconductor—Typical
examples of pentavalent impurities are arsenic
(As, Z = 33) and antimony (Sb, Z = 51). In the
crystal, they will be surrounded by the atoms of
the semi-conductor material. Four (out of five)
valence electrons of the neighbouring Ge atoms.
The fifth valence electron of impurity (say, As)
finds no place in covalent bands and remains to
move randomly in the crystal lattice. Thus, each
impurity atom donates a free electron to the
semiconductors. Hence, the impurity is called the
donor impurity. The semiconductor containing do-
nor type impurity is called n-type semiconductor
because it has negative charge carrier (electrons).
Ge Ge Ge
Free
electrons
Ge As Ge
Ge Ge Ge
Pentavalent
Impurity atom
In n-type semiconductors, electrons are
majority charge carriers but there are still present
some thermally generated holes which are called
the minority charge carriers.
(ii) p-type semiconductors—If a small
amount of trivalent impurity (having three valence
electrons) is added to a pure semiconductor it is
known as p-type semiconductor.
In extrinsic of current flow—Typical
examples of trivalent impurities are gallium (Ga,
atomic number 31) and indium (In atomic number
49). When a small amount of such impurity is
added to germanium crystal, three valence elec-
trons of impurity atoms form covalent bonds with
three valence electrons from three neighbouring
germanium atoms. There is a deficiency of one
636 | Physical Sciences
electron to complete the fourth bond. This electron
deficiency is called the hole and it behaves like a
positively charged particle since there is a strong
tendency of semiconductor crystal to form cova-
lent bonds, therefore, a hole attracts one electron
from a nearby covalent bond.
Ge Ge Ge
Hole
Ge In Ge
Ge Ge Ge
Trivalent
impurity atom
Such an impurity that produces holes inside a
semiconductor crystal is called acceptor impurity
because it accepts electron from the nearby
germanium bond. The semiconductor containing
acceptor type of impurity is called p-type semi-
conductor because the charge carriers are positive
holes.
It should be noted that p-type and n-type
germanium are both electrically neutral because
they are produced by adding neutral impurity
atoms to neutral germanium or silicon.
Fermi Level in Semiconductors
(a) In intrinsic semiconductors—The energy
band diagram of an intrinsic semiconductor is
shown below—
In above figure Ec and Ev represent respec-
tively the lowest energy in C.B. and highest
energy level in V.B, E g is the bond gap energy.
The other diagram above shows the Fermi-Dirac
distribution. At 0 K, the probability of finding an
electron in the C.B. and hole in V.B. is zero.
At room temperature (≈ 300K) some of the
electrons fro m the V.B. move to C.B. and occupy
the states near the bottom of C.B. This gives the
non-zero value of distribution function f (E), i.e.,
the tail of the distribution curve extends into the
C.B. Similarly, an equal number of holes are
present near the top of the V.B. and the probability
of occupancy decreases from unity. Since, Fermi
level can be defined as that energy level at which
the probability of occupancy becomes half. We
find that the Fermi level for intrinsic semicon-
ductor is near the middle of the forbidden gap.
(b) In extrinsic semiconductor—The band
diagram for an n-type semiconductor is shown in
figures below, where ED represents the energy
level corresponding to the donor.
In case of an n-type semiconductor, if we take
that all donor atoms are ionized, then the donor
electrons will occupy the states near the bottom in
C.B. Then a smaller number of electrons from
V.B. will be able to cross the band gap by thermal
agitation. Hence, the number of holes in V.B. is
decreased. Since, the Fermi level gives the energy
state for which the probability of occupancy is
half, E f for an n-type semiconductor is close to
C.B. Thus in n-type semiconductor Fermi level
will move towards the C.B. If the doping is very
large the Fermi level may merge with the C.B.
The same kind of reasoning when applied to a
p-type conductor, leads to the result for a p-type
semiconductor, the position of the Fermi level is
located near the V.B.
In a p-type semiconductor, semiconductor
Fermi level will move towards the valence band
and for sufficiently large doping, it may move into
V.B.

When there is increase in temperature of the
semiconductor, then both n-type and p-type semi-
conductor become intrinsic, we may conclude that
with temperature rise, the Fermi level EF moves
towards the centre of the forbidden gap for both
types of semiconductors.
Conduction of current in semiconductors—
In a semiconductor each carrier (electron or hole)
moves in a random way owing to its thermal
energy when no external field is applied to semi-
conductors. The frequent change in direction of
the path of the carrier occurs due to Coulombs
field of the ionized donor and acceptor atoms and
scattering by the vibration of the lattice atom.
When an external field is applied then this field
superimposes on the random motion of the
carriers. In the steady state, the rate of momentum
gained from the field equals the rate of loss of
momentum due to scattering. Hence a steady drift
velocity giving rise to a steady flow of current.
(a) Current flow in a biased n-type semi-
conductor—Suppose a voltage be applied to an n-
type semiconductor placed between a pair of elec-
trodes. Now there will be a steady drift of the free
electrons towards the positive electrode due to the
field produced by the voltage. The electrons reach-
ing the positive electrodes disappear at the metal
of the electrode and the immobile positive ions in
the neighbourhood of the negative electrode
remain unneutralised due to the dirft of free
electrons. These ions immediately attract electrons
from the negative electrode. Hence a continuous
flow of electrons from one terminal of the voltage
source to the other terminal via the semi-
conductor occurs.
The rate of flow of electrons from the
negative electrode to the semiconductor and from
the semiconductor to the positive electrode is
determined by the applied voltage and the
conductivity of the semiconductor which account
for current flow.
(b) Current flow in a biased p-type semi-
conductor—When a voltage is applied across a
Physical Sciences | 637
pair of electrode and a p-type semiconductor is
placed between these electrodes, a constant drift
of the hole which are here majority carrier; takes
place towards the negative electrode.
On reaching the negative electrode they com-
bine with the electrons arising out of the metal of
the negative electrode and disappear. Consequently
an equal number of holes is generated near the
positive electrode. As the holes drift away from
the positive electrode, they leave behind the
unneutralized immobile negative charges. These
charges and the positive electrode give rise to an
electric field which causes the ionised electrons to
leave the acceptor atoms and draws them towards
the positive electrode where they are lost. By
loosing an electron in this process the acceptor
atom tries to clinch an electron from an adjacent
band hence a hole is formed. The formation of
holes takes place at the rate equal to the rate at
which they disappear at the negative electrode.
This rate is controlled by the magnitude of the
applied voltage and the conductivity of the semi-
conductor.
Mobility
It may be defined as the average drift velocity
per unit electric field. If ‘v’ is the drift velocity of
a carrier in an electric field of ε volt/m, then
mobility (µ) may be given by
v
µ = = m2sec–1 volt–1
ε
Here the charge carriers experience forces not
only due to this applied field but also due to field
produced by internal crystal potential which
modifies the mass of carriers in such a way that it
obeys the laws of mechanics.
Now, according to Newton’s law when a field
ε is applied, the rate of gain of momentum by the
charge e is eε and the rate of loss of momentum in
m*v
scattering processes is given by where
τ
m* = effective mass, v = drift velocity of carriers
and τ = relaxation time. In the steady state
638 | Physical Sciences

m*v eτ also along positive x-direction. The flow gives rise

eε = or µ =
τ m* to diffusion current as given by :
e〈τ〉 dp
Here it can be shown that µ = where Jp = – eDp
m* dx
〈 τ 〉 shows the average value of τ involving Here e is electric charge, Dp is diffusion co-
distribution function of carriers. It is also found efficient. A similar expression can be written for
that µp < µ n ; it means that holes are less mobile diffusion electron current.
than the electrons. In presence of both electric field and con-
centration gradient, the hole current is given by
Current Density
Jp = pe µp ε – eDp dp/dx
Current per cross-sectional area or the charge
crossing unit cross-sectional area per second and electron current is
conductivity is called current density (J). In other dn
Jn = ne µn ε + eD n
words, conductivity may be defined as the current dx
density per unit applied electric field. If J is the
current density due to an applied electric field ε, Relation between Diffusion Current
the conductivity will be : and Mobility
σ =
J
or J = σε The quantities D and µ are related by the
ε following equation, known as Einstein relation
Current density may also be explained by Dp Dn
giving =
µp µn
J = nev
kT
where n is expressed in m–3 , v in ms–1 and e in =
coulomb. e
So, J = nev = ne µε where k is Boltzmann constant and T is absolute
J ne µε temperature. At room temperature :
Thus, σ = = = ne µ
ε ε Dp D
= n
Overall conductivity of a semiconductor µp µn
containing electrons and holes is given by = 0·026V
σ = σn + σp Hall Effect
= en µn + ep µp When a piece of conductor carrying a current
= e (n µn + p µp ) is placed in a transverse magnetic field, an elec-
tric field is produced inside the conductor in a
Diffusion
A particle in the presence of a concentration
gradient tends to move from a region of higher
concentration to region of lower concentration in
order to destroy the non-uniformity in concentra-
tion. This movement of particles is called the
phenomenon of diffusion.
Suppose there be concentration of holes vary-
dp direction normal to both the current and the
ing with distance x in the semiconductor, being magnetic field. This phenomenon is called Hall
dx
dp effect.
the concentration gradient. Suppose is negative
dx Let us consider a rectangular piece of
it means that the hole concentration is greater for conductor carrying a current I in the positive X-
lower values of x. The holes move at random due direction and a magnetic flux density B in the
to their thermal energy but there is diffusive flow positive Z-direction.

In the negative Y-direction the current carriers
will experience a Lorentz force. In the metal or n-
type semiconductor the current carriers are elec-
trons so the electrons will accumulate at the
bottom surface making it negatively charged with
respect to top surface. Thus, an electric field
called the Hall field will be developed along the
negative Y-direction. The force on the current
carrying electrons due to the Hall field opposes
Lorentz force and at equilibrium they balance each
other, while for p-type semiconductor the current
carriers are holes hence the accumulation of
carriers on the bottom of surface make the surface
positive relative to the top surface.
Therefore, the Hall field is produced along the
positive Y-direction the force on the holes due to
Hall field opposes the Lorentz force and balances
in equilibrium.
If ε H = Hall field, v = Drift velocity in X-
direction.
Then Hall force on carrier of charge
e = eεH
Lorentz force on that carries = Eev
Current denisty in X-direction J = nev
BJ
Thus, εH = = vB
ne
The Hall effect is described by means of Hall
coefficient RH, that is
εH = R HJB
Hence, we obtain
ε BJ 1
RH = H = = Let us consider two different metals A and B
JB neJB ne
For electrons
1 kinetic theory, the electronic pressure in A is
RH = –
ne
where n gives the concentration of electrons for
holes
1 thereby creating a difference of potential at the
RH =
pe
Here p gives the concentration of hole. Thus,
sign of RH tell whether the sample is an n-type or
p-type semiconductor.
Calculation of RH
R H is calculated by measuring the Hall
voltage which generates the Hall field. If VH is the
Hall voltage across a sample of thickness d, then
VH = εH d = RH J.B.d
Physical Sciences | 639
If w = Width of the sample then its cross-
sectional area = w.d
I ,
Thus, J =
wd
R IB
So, VH = H
w
VHw
or RH =
IB
V H will have opposite signs for n-type and
p-type semiconductor.
Importance of Hall Effect
The Hall effect measurements provide the
following information about the solid—
1. The sign (electrons or holes) of charge
carrier is determined.
2. The mobility of charge carriers is measured
directly.
3. The carrier concentration (number of
charge carriers per unit volume) can be calculated.
4. It can be used to determine whether the
given material is metal, insulator or semi-
conductor.
5. From the knowledge of measured Hall
voltage, the unknown magnetic field can be
measured provided the value of Hall constant for
the slab is known.
Thermoelectric Effects
On the assumption that the electron density
(number of free electrons per unit volume) varies
from metal to metal; the electron theory is able to
explain thermoelectric phenomenon also.
in contact and the electron density in A be greater
than that in B. So that according to principles of
greater than that in B. Due to this difference in
electronic pressure electrons will diffuse from A
to B, making A positive and B negative and
junction of the two metals. This potential
difference will tend to prevent further electrons
passing from A to B; when it reaches a certain
value, a state of equilibrium sets in and the migra-
tion of electrons ceases completely. Thus, we get a
simple explanation of the potential difference
created at the junction of two metals. This natural
potential difference is physically distinct from the
so called contact potential difference of two
metals. This general conclusion can be applied to
the two effects of thermoelectric phenomenons.
640 | Physical Sciences
Peltier effect comes into existence when an
external potential difference is applied to the
junction, the current flows from A to B. Since the
electron density in A is assumed to be greater than
that in B, the movement of electrons of electrons
from B to A due to the passage of current involves
the work to be done against the difference in
electronic pressure and hence the energy will be
absorbed at the junction which in consequence
gets cooled. If the direction of the current is
reversed, the electrons moves from A to B which
makes energy available at the junction in the form
of heat, i.e., the junction gets heated.
Thomson effect refers to the absorption or
evolution of heat due to the passage of a current in
a single unequally heated conductor. It is some-
what complicated on account of variation of its
sign with different metals copper, cadmium, silver,
zinc are said to have a positive Thomson effect;
since in these metals heat is evolved when current
flows from hot to cold side while heat is absorbed
when the current flows from cold to hot side. On
the contrary, iron, cobalt, bismuth and nickel are
said to have a negative Thomson effect; since the
phenomenon is reversed in them. In lead, the
Thomson effect is probably nil.
Dielectric
A dielectric is a material which contains no
free electrons, so that no current can flow through
it. As a result the electrical conductivity of a
dielectric is poor and for an ideal dielectric, it is
zero. According to band theory of solids a
dielectric is a material in which the energy band
gap between valence and conduction band is more
than 3eV. Glass, plastic, mica, rubber, wood,
turpentine oil etc., are examples of dielectrics.
Dielectric Polarisation
When a polar or a non-polar molecule is
placed in an external electric field, the small dis-
placements of the orbital electrons will cause the
distance between the centres of gravity to alter.
Thus, non-polar molecules become induced
dipoles whereas polar molecules, which are
already dipoles, will be oriented by the field and
may have resultant dipole moment modified
since they possess—(i) Permanent dipole moment,
(ii) Induced dipole moment. The orientation of the
induced dipoles or of the permanent dipoles in an
external electric field will be such as to set or tend
to set the axis of dipole along the lines of force.
This phenomenon is called electric polarisation.
The behaviour of a dielectric in the field between
two parallel plates is shown in given fig.
+ –
-+ -+ -+
-+ -+ -+
-+ -+ -+
-+ -+ -+
(a) In the absence (b) In the presence
of the field of the field
When no electric field is applied the molecules
do not show any polarity since they are unelectri-
fied non-polar molecules or they will be polar
molecules which will be oriented at random and
will be in such large numbers that the resultant
electric effect due to all of them is zero. When a
potential difference is applied across the parallel
plates, the electric field is produced. This electric
field tends to align the induced or permanent
dipoles along the electric lines of force, the degree
of alignment depending on the field strength. The
negative charge on the dipoles will face the posi-
tively charged plate, while the positive charges
will face the negative plate. The induced field in
the dielectric is, therefore, opposite in direction to
the external applied field.
Let us now consider a slab of dielectric,
placed between plates such that there is a gap
between the plates and the walls of the dielectric.
Such a gap is always present, through it may be
only of molecular dimensions in width, even
through the dielectric is said to fill up the space
between the plates. If the potential difference is
applied between the plates, an electric field is
produced. This electric field produces an induced
field within dielectric in direction opposite to the
external field. Hence, the field in the dielectric is

less than the original field in space. The induced
charges within the polarised dielectric are called
fictitious or bond charges whereas the charge on
the plates are called real charges. If the dielectric is
homogeneous and isotropic, then all the molecules
are polarised to the same extent. In this case the
fictitious charges within the main body of the
medium is adjacent to the positive side of its
neighbour. However, at the surface of the
dielectric in contact with the plates, the fictitious
or bound charge are not neturalised; thus electric
field induces charges only on the surface of the
dielectric. These induced charges are opposite to
the charges on the plates; consequently the field
within the dielectric is reduced as compared with
that in free space.
Magnetism
A charge separation constitutes an electric
dipole and the polarisation represents the average
density (per unit volume) of the electric dipole
moment in a given direction while the electric
polarisation has to do with the static positioning of
the charges, the magnetisation is due to the motion
of these charges. A plane current loop is equiva-
lent to a magnetic dipole so far as the magnetic
effects both produced and experienced by them,
are concerned. Electrons in an atom move in
certain prescribed directions about the nucleus of
the atom. The electronic motion can be described
by assigning a certain angular momentum to each
electron. Consequently the magnetic properties of
a substance are determined by the state of angular
momentum which the electrons may possess.
Classification of Magnetic Materials
The magnetic properties of a substance
depend upon—
(a) Whether the substance has permanent
dipoles.
(b) If it has, how these dipoles are oriented
with respect to each other. Based on these criteria,
the magnetic materials are generally classified into
the following five categories :
1. Diamagnetic
2. Paramagnetic
3. Ferromagnetic
4. Antiferromagnetic
5. Ferrimagnetic
Physical Sciences | 641
Importance in Solid State Physics
Many magnetic materials contain more than
one kind of magnetic species because magnetic
ions located at different crystal sites may behave
differently, according to their positions in the
lattice. The interaction between neighbouring
magnetic ions thus becomes an important para-
meter for the studies of internal field and lattice.
Diamagnetism
Classically, the diamagnetism i s related to
changes in the orbital motion of electrons known
as Larmor precession that occur when atomic
systems are placed in a magnetic field. Atoms or
ions in which the orbital and spin angular
momenta cancel in pairs have angular momentum
J = 0 and no permanent magnetic moment to cause
paramagnetism.
Diamagnetism is associated with the tendency
of electrical charges partially to shield the interior
of a body from an applied magnetic field. It may
be recalled that the current induced in a closed
electrical circuit by a magnetic field is always in
such a direction as to tend to keep the total flux
unchanged. Thus, the circuit has in affect a
negative susceptibility. The effect is retained even
in systems of charges that must be treated by
quantum mechanics and is responsible for
diamagnetism.
Since, all atoms or ions produce a diamag-
netic contributions to the total susceptibility,
although it may be masked by the other types; it is
consequence of the magnetic moment induced in
the atoms by an external field. In this respect
diamag-netism may be compared with the
electronic polarization in an electric field. Both
are essentially independent of temperatures. There
exists however an essential difference; in the
electrical case the induced moment lies along
the direction of an applied field leading to a
positive electrical susceptibility, in the magnetic
case the induced moment produces a negative
susceptibility. It should however be noted that the
diamagnetic susceptibility is determined essen-
tially by charge distribution in atoms.
Paramagnetism
Paramagnetism is related to the tendency of a
permanent magnet to align itself in the direction of
magnetic field such that its dipole moment is
parallel to the field. In atomic systems the perma-
nent magnetic moment is associated with the
642 | Physical Sciences
electronic spin as in atoms, molecules and lattice
defects possessing an odd number of electrons, the
total spin cannot be zero. Hence, a paramagnetic
substance has a non-vanishing angular momen-
tum. It may also be the permanent moment of
unfilled atomic shells that arises from a combina-
tion of spin and orbital momentum.
Cause of Paramagnetism
Some atoms and ions have permanent mag-
netic moments. In the absence of magnetic field,
these moments usually point in random direction
and hence producing no macroscopic magneti-
sation. But in the presence of magnetic field, the
moments tend to line up preferentially in the field
direction and produce a net magnetisation. Since,
the moments line up in the direction of the field
and enhance the external field, the susceptibility is
greater than zero. The paramagnetism depends
upon magnetic moment of ions, the origin of these
magnetic moments is important. There are two
chief origin for the moments, one from the spin
motion and the other from the orbital motion of
the electron.
The orbital moment arises in principle because
an electron in an orbit about an atom can be
considered to be a small circulating current about
the nucleus. The magnetic moment is related to
the angular momentum by an important relation as
follows—
The primitive magnetic moment µm of a
circulating loop is given by
µm = iA.
where i the current in the loop and A is the area
enclosed by the circulating current.
Classical Theory of Paramagnetism
Langevin considered a paramagnetic gas in
which each particle is associated with a permanent
magnetic moment µ. Hence, when a mass of gas
was placed in a field, each particle tended to
set itself along the magnetic field. Now, at the
ordinary temperature the particles are subjected
to ordinary thermal agitation so the particles
would be prevented from taking exact alignment
but a kind of statistical equilibrium will be set up.
Thus, the Maxwell Boltzmann statistics was
applied. Then potential energy of a particle
making an angle θ with the field H was calculated
from this magnetic moment and then susceptibility
was calculated.
Drawbacks of Classical Theory
(1) If the permanent magnetic moment is
associated with and is proportional to the angular
momentum of the moving electronic charges, then
the magnetic moment cannot have a fixed magni-
tude but must take all values ranging from – ∞ to
+ ∞.
(2) If the particle of magnetic moment µ is
placed in a weak magnetic field H, it cannot place
itself at any angle θ to the field as was supposed
by Langevin.
Weiss Theory of Paramagnetism
Weiss assumed that since in a real gas the
molecules are mutually influenced by their
magnetic moments, there should exist within the
gas a molecular field, produced at any point by all
the molecules in the neighbourhood proportional
to and acting in the same sense as the intensity of
magnetisation. If this molecular field be Hm, then
H m = AI, where A is molecular field coefficient
and I is the intensity of magnetisation.
According to Weiss theory,
cm
χm =
T–θ
This is the Curie-Weiss law. θ is the Weiss
constant; also known as Curie temperature at
which ferromagnetics become paramagnetics.
Quantum Theory of Paramagnetism
According to quantum theory the permanent
magnetic moment of a given atom or ion is not
freely rotating but restricted to a finite set of
orientation and relative to the applied field. In
quantum theory the possible value of resolved
magnetic moment in any direction is given by
MJgµB .
where MJ = J, J – 1, J – 2, … J
Magnetic moment parallel to H = MJgµB H
Potential energy = – MJgµB H
and the Magnetic susceptibility is given by :
N2 µ J 2
χm =
3RT
This expression is identical with the classical
result
N2 µ 2
χm =
3RT
 Physical Sciences | 643

Here the value of µ has been replaced by µJ.
This relation of µJ2 is also in agreement with other
quantum relation. This expression for χ deduced
on the basis of quantum theory gives the value of
effective Bohr magneton number
absence of an applied magnetic field. The Curie
point is the temperature at and above which the
spontaneous magnetisation vanishes Weiss
explained the theory of ferromagnetism on the
basis of two hypothesis :

µB .pB = µ = 
√
(3kT χ) (a) A ferromagnetic specimen contains in
general a large number of small regions called
or pB = g √ (

J (J + 1) ) domains which are spontaneously magnetised.
and Weiss magneton number This magnetisation of the specimen is determined
by the vector sum of the magnetic moments of the
pw = 4·97 pB individual domain.
= 4·97g 
√
J (J + 1) (b) Within each domain the spontaneous
This is known as Hund’s formula. magnetisation is due to existence of molecular
Now two cases arises— field which tends to produce a parallel alignment
(a) When the energy of transition between of the atomic dipoles.
two states is very large as compared to kT, i.e., Existence of magnetic moment leads to
‘wide multiplets’. cooperation or interaction between the atomic
(b) Energy of transition is very small as dipoles. Tendency of cooperation is to produce
compared to kT, i.e., ‘Narrow multiplets’. parallel alignment. After a long derivation, we
find the following result—
Ferromagnetism
JgµB H
Ferromagnetic substances have following B(a) = …(A)
kT
properties—
where B(a) is the Brillonin function, J is angular
1. They acquire a relatively high magnetisa- momentum, µB is Bohr magneton.
tion in weak fields.
2. They do not exhibit a linear proportionality In the case of ferromagnetism H is replaced
between the magnetisation and the fixed strength. by
3. They can exist in a permanently magnetised Heff = H + AI
state even when there is no external field. JgµB Heff
So, B(a) =
4. Their magnetisation is not a unique function kT
of the field strength but depends also on the fields
JgµB (H + AI)
to which they have been subjected previously. =
kT
5. They lose their special characteristics above
a certain critical temperature and then behave as For spontaneous magnetisation H = 0
paramagnetic. JgµB AI
6. If the magnetising field is small, the ∴ B(a) =
kT
susceptibility decreases as the low temperature
are approached; but if the field is large, it
increases.
or I =
B(a)kT
JgµB A [
= M ˙.˙ I =
M
V ]
…(B)

Therefore, ferromagnetism is probably to be Equation (B) represents a straight line; the

regarded as a special manifestation of paramag-
netism limited to a range of temperature below
Curie point.
Quantum Theory of Ferromagnetism
Ferromagnetic substances are those sub-
stances which possess a spontaneous magnetic
moment, i.e., a magnetic moment even in the
slope of the line is proportional to temperature.
From given figure it follows that for T < θf, one
obtains a non-vanishing M value although H = 0.
For T = θ f the slope of the line represented by
equation (B) is equal to that of the tangent of
curve (2) at the origin. Thus, for T ≤ θ f the two
curves do not intersect and hence there is no
spontaneous magnetisation.
644 | Physical Sciences

that the A spins tended to line antiparallel to the B

spins. At low temperatures this interaction is very
effective and in an external field the resulting
magnetisation will be small. As the temperature is
raised, the efficiency of the interaction becomes
less pronounced and the susceptibility increases.
Finally a critical temperature TN, known as Neel
temperature, is reached above which the spins are
free and above this temperature the antiferromag-
Schematic representation of the method for finding netic material becomes paramagnetic.
the spontaneous magnetisation of temperature T. The direct evidence for the basic picture of
antiferromagnetism has been obtained from neu-
Antiferromagnetism tron diffraction experiments.
As we know that the behaviour of a ferro-
magnetic material may be explained quantum
mechanically in terms of the exchange interaction.
It was however assumed in this theory that the
exchange integral J is positive. When the exchange
integral J is negative favouring an ordered anti-
parallel orientation of electron spins, one has an
antiferromagnetic substance.
↑↑↑↑↑↑↑↑ ↑↓↑↓↑↓↑↓↑↓ (a) Curie law
Exchange Integral > 0 Exchange Integral < 0
Ferromagnetism Antiferromagnetism
(a) (b)
Such systems were first investigated theoreti-
cally by Neel and Bittei. Experimentally antiferro-
magnetism was first discovered as a property of
MnO by Bizette, Squire and Tsai in 1938.
The most characteristic property of antiferro- C
magnetism is that its susceptibility shows a maxi- χ =
T – TC
mum as function of temperature as shown in given (b) Curie Weiss law
figure. (T > T C)
χ
ANTIFERROMAGNETISM

−θ O TC T
The molar susceptibility χM of MnF 2 as function of C
χ =
temperature. T+θ
C
This feature may be explained on the basis of χ = (T > TC)
T+θ
following model. Let us consider a crystal (c)
containing two types of atoms A and B. Let the A
atoms occupy the corner points and the B atoms Distinguishing features of temperature depen-
being located at the centres of these cubes. Further- dence of magnetic susceptibility in para, ferro and
more, let the interaction between the atoms be such antiferromagnetism.

Ferrimagnetism
Ferrites form a class of solid oxides of iron.
One way of writing their formula is Me++ Fe2+++O4.
Me++ stands for any suitable divalent metal
ions such as Co, Mn, Ni, Cd, Mg, Ni for Me gives
Nickel Ferrite, Fe for Me, Fe++, Fe2 +++O4 gives
Ferrous Ferrite commonly known as Fe3 O4 and its
natural mineral is magnetite, which was originally
known as loadstone.
Many ferrites have mangetic property strongly
resembling with those of metallic ferromagnetics.
They show hysteresis, spontaneous magnetisation
and associated phenomenon and their saturation
momenta are of the same order as iron. Their
technical interest is because of the fact that they
are not metals but ionic crystals which put them in
the insulator or semiconductor class. This means
that power loses due to eddy currents are much
reduced and ferrites are, therefore, well suited to
high frequency application.
Superconductivity
Kammerlingh Onnes found that the electrical
resistance of some metals, alloys and compounds
drops suddenly to zero when the specimen is
cooled below a certain temperature. This phe-
nomenon is called ‘superconductivity’ and the
cooled specimen is called a ‘superconductor’.
The critical temperature TC below which a
material undergoes a transition from a state of
normal conductivity to a superconducting state is
different for different materials. It varies from
23·2K for the alloy Nb3 Ge to 0·01K for some
16 TIN
(Superconductor)
12
Resistivity
8
4 Silver
(Normal Metal)
0
TC 4 8 12 16 20 T
Temperature (K)
semiconductors. The normally good conductors
like Cu, Ag, Au, Li, Na, K etc., do not show
super- conductivity even at temperature as low as
a small fraction of 1K. These metals are called
normal metals. Figure shows the resistivity at very
low temperatures for a superconductor tin (TC =
3·7K); and a normal metal silver.
Physical Sciences | 645
Type of Superconductors
The superconductors can be grouped into two
distinct classes according to their behaviour in an
external magnetic field.
1. Soft superconductors or type–I—The
superconductors in which the magnetic field is
totally excluded from the interior of superconduc-
tors below a certain magnetising field H C and at
HC the material loses superconductivity are termed
as soft or type-I superconductors.
Examples—The examples of type–I super–
conductors are tin, lead, mercury.
Important Features of Soft Supercon-
ductors—
(i) They exhibit complete Meissner effect.
(ii) The critical value of magnetic field HC at
which magnetisation drops, are very low. The
maximum known critical H C for type-I
superconductors are of the order of 1000 gauss.
Low value of HC makes these materials unsuitable
for use in high field superconducting magnets.
The magnetisation curve shows that if the
magnetic field is reduced below HC , the material
again acquires superconducting property and the
field is expelled.
2. Hard supercondcutors type-II—The
superconductors in which the material loses
magnetisation gradually rather than suddenly, are
termed as hard or type-II superconductors. Its
important properties are—
(i) They do not show complete Meissner
effect.
(ii) These superconductors have two critical
fields HC1 and H C 2 called lower and upper critical
fields respectively.
The specimen is diamagnetic below HC 1 . At
HC1 the flux begins to penetrate and the penetration
of flux increases until the upper critical field HC2 is
reached. At H C 2 the magnetisation vanishes and
specimen returns to normal conducting state. The
646 | Physical Sciences

field H C 1 occurs before then HC . The value of →

Hence, if the resistivity goes to zero while j
critical field HC2 for type-II materials may be 100
→
times or more higher than the value HC for type-I is held finite then E must be zero. This means
superconductors. →
E = 0, inside a superconductor. Then, we get
Examples—In homogeneous material and
→
include Nb-Zr, Nb-TC alloys and Va-Ga and Nb- ∂B
= 0
Sn intermetallic compounds. Type-II super- ∂t
conductor are technically very useful than type-I
→
superconductors due to tolerating high magnetic or B = Constant
fields. The magnetisation curve for these super-
which shows that the flux through the metal
conductors is given in figure.
cannot change on cooling through the transition.
But Meissner showed that a change does occur as
the material passes over to the superconducting
state. The Meissner effect shows that a super-
conductor behaves as if inside the specimen we
→
have B = 0. Hence, it suggests that perfect
diamagnetism is an essential property of the
superconducting state.
Therefore, the conditions defining the super-
The Meissner Effect conducting state are—
Meissner discovered that if a superconductor →
E = 0 (From the absence of resistivity)
is cooled in a magnetic field down to the transition
temperature TC; then at the transition the lines of →
B = 0 (From Meissner effect)
induction B are pushed out. This phenomenon is
Now at
called Meissner effect. It shows that in an applied
external field, a bulk superconductor behave as if →
B = 0
inside the specimen B →
From the relation 
H H
H′
→ → → 
 B = H + 4πmv 
mv = –
Meissner 4π
effect
AA where H′ is the actual field at the surface of the
specimen and mv is the magnetic moment per unit
volume. Here –ve sign indicates that a supercon-
T > TC or H > HC ductor exhibits a perfect diamagnetism.
T < TC or H < HC
(a) (b) (b) Explanation for the persistent currents
(a) Meissner effect contradicts the set up inside the superconductor according to
Maxwell’s equations—Before the Meissner’s Meissner effect—Due to the induction of current
discovery the Maxwell equations were thought to in the specimen, a flux at the transition temperature
be quite adequate to explain the phenomenon of is ejected which produces a flux equal and oppo-
superconductivity. From Maxwell’s equations, we site to the original flux due to the magnetic field.
get Since, in the superconducting state resistivity is
→ → → – ∂B zero, the current due to induction in the specimen
curl E = ∇ × E =
∂t persists as long as the field is on. The specimen
and from Ohm’s law acquires a negative moment producing a field in a
 → direction opposite to the direction of applied field
→ → → → E due to persistent current. When the external field
E = e j j =σ E = 
 e is removed, persistent current goes off.

From the result of above experiment it seems
that the persistent current can be set up only in the
presence of applied field. Meissners showed that
persistent currents can be set up in a sample of
suitable shape. Let the specimen is in the shape of
an anchor ring and is placed in a magnetic field.
In figure,
(a) The specimen is placed in a magnetic
field. The flux will penetrate all space.
(b) Now the temperature of ring is reduced
down to the transition temperature. Therefore,
specimen due to the occurrence of superconduct-
ing state, will eject the flux and the new flux
distribution will be shown in fig. (b).
(c) This explains that irrespective of the
presence or absence of applied field, persistent
currents are maintained in the superconductor.
(d) On removing magnetic field no current
flows.
Since, the magnetic induction inside a
specimen is zero, it means the change from the
superconducting to the normal state or normal state
to the superconducting state by the application or
removal of a field ≥ HC is a reversible process.
Applications
There are many uses of superconductivity.
1. Superconducting cables can be used to
transmit electric power over long distances without
resistive losses. This would be economical, if the
cost of keeping the cable below its critical tem-
Physical Sciences | 647
perature TC is less than the value of the power that
would be lost using ordinary cables.
2. Since, superconductors are diamagnetic in
nature, they can be used to shield out unwanted
magnetic flux, as in shaping the magnetic lens
system of an electron microscopes.
3. Superconducting electromagnets carry
large resistanceless currents and hence produce
large magnetic fields, can be constructed if we use
superconducting wire for magnet windings,
currents of the order of 100 amperes can be made.
B. C. S. Theory of Superconductivity
In 1957 Basdeen, Cooper and Schriffer (BCS)
proposed a theory of superconductivity in which
they included the electron-phonon and electron-
electron interactions. The basic features of this
theory can be given by following way—
An electron in a solid passing by positive ions
in the lattice exerts upon them coulomb attractions
which give them momentum and cause them to
slightly together. The region of increased positive
charge density propagates through the lattice as an
elastic wave carrying momentum supplied by the
electron. Now, when a second electron is passing
by the moving region of increased positive-charge
density, it will experience an attractive coulomb
interaction and absorb all the momentum carried
by the moving region. In this process the net effect
is that the two electrons have exchanged some
momentum with each other. Hence, they have
interacted with each other through the lattice
deformation. This electron-electron interaction is
attractive because both the electrons participated
in coulomb attractions with the lattice. If the
temperature is low enough and the two electrons
have antiparallel spins, the attraction between the
electrons exceeds slightly the direct coulomb
repulsion between them. Then the electrons are
weakly bound together and form a so-called
‘cooper pair’. At low temperature a large number
of such pairs are formed, if there are enough
conduction electrons lying just below the Fermi
energy level, when an external electric field is
applied, the cooper pairs move through the lattice
maintaining their order, as if the motion of each
pair is locked into the motion of all the rest.
Therefore, electrons move in highly correlated
pairs, without involving in the random scattering
from lattice imperfections that cause electrical
resistance. Thus, the solid is a superconductor.
648 | Physical Sciences
Superconductors have a relatively strong
electron-lattice-electron interaction. On the other
hand, metals can be good conductor when their
conduction electrons have a weak interaction with
the lattice. That is why the metals which are
normally good conductors do not become super-
conductors even at very low temperature.
Explanation for the Attractive Inter-
action
The electrons and ions move very quickly.
The electrons are always passing through the field
of ions which are themselves vibrating rapidly and
the effect of the virtual phonon on the ionic
vibration depends very much on their relative
frequencies. Now, two cases arises—
Case I—Let us denote lattice frequency by ν l
and the virtual phonon frequency by νp.
If νl ≥ νp, the system will tend to oscillate so
that it is in phase with the driving phonon. Thus,
the original charge fluctuation which occurred
when electron l changed its state will be over
compensated by the ionic motion when the ions
move in phase so that the coulomb repulsion on
electron 2 is diminished.
Case II—If Vl ≤ Vp —The ionic fluctuation
will be out of phase with the original virtual
phonon. Hence, the original charge fluctuation
which occurred when electron, changed its state
will be under compensated by the ionic motion
when the ions move out of phase. So, that
coulomb repulsion on electron 2 is magnified.
Under the first case it is possible for the ionic
influence to overcome that of the original electron
charge fluctuation so that electron 2 is actually
attracted to electron 1 rather than being repelled
from it. If this occurs, the energy of the system
will be reduced.
Conditions of Attraction to Take Place
If the interaction between electrons and
phonons was very weak, then there would be little
likehood of a virtual phonon exchange. The subs-
tances such as the alkali and interact strongly with
the phonons, are not superconductors whereas ma-
ny of the bad conductor are the superconductors.
Penetration Effects
The expulsion of an applied magnetic field
from a superconductor can be considered to be
due to the opposing effect of the field produced by
currents which flow on the surface of the speci-
men. It is clear however, that when we consider
some type of surface current, the current does not
flow only on the surfaces—there must be some
small but finite layer in which the current flows.
Only within this layer will the full influence of the
magnetic field due to the surface currents be felt
and a proper Meissner effect be produced. In the
layer itself the magnetic induction will not be
zero. Hence, we might accept that an applied
magnetic field penetrates into the specimen a short
distance. Such a penetration which has been
observed was predicted by a theory of supercon-
ductivity proposed by F and H London. This
theory essentially suggest relations between the
superconducting current and the electric and
magnetic fields associated with that current.
Explanation for the Behaviour of
Superconductors
The behaviour of a superconductor as a
perfect diamagnet in an external magnetic field
can be explained. According to Lenz’s law, when
the magnetic flux through a circuit is changing, an
induced current is established in such a direction
as to oppose the change in flux. The external
magnetic field does not penetrate the interior
of the superconductor because in it the conduction
electrons, moving in cooper pairs without any
resistance, adjust their motion to produce a
counteracting magnetic field. Obviously, it is
necessary for a superconductor to have a
persisting current to keep the magnetic flux out
when there is an external magnetic field. Hence,
the two main properties of superconductors
namely the exclusion of magnetic flux and the
absence of resistance to current flow, are related
to each other.
Again when the external magnetic field is
increased beyond critical value, the superconductor
becomes normal. It can also understood when the
magnetic field is turned on, the superconductor
acts to exclude this field from its interior. The
energy decrease of the magnetic field appears as
increased energy of the electrons in cooper pairs.
As the magnetic field is increased, the energy
acquired by the superconductor also increases.
Beyond the critical value of the field, the energy
of the superconducting state becomes higher than
the energy of the normal state and the material
becomes normal. In this way nature of super-
conducing materials have been explained.
 Unit–10

NUCLEAR AND PARTICLE PHYSICS

General Properties of Nucleus
Nuclear Size
Rutherford’s work on the scattering of α
particles showed that the mean radius of an atomic
nucleus is of the order of 10–14 to 10–15 m while
that of the atom is about 10– 1 0 m. Thus, the
nucleus is about 10,000 times smaller in radius
than the atom.
The empirical formula for the nuclear radius
is
R = r0 A1/3
where A is the mass number and r0 = 1·3 × 10–15 m
= 1·3 fm. Nuclei are so small that the Fermi (fm)
is an appropriate unit of length. If 1 fm = 10–15 m,
from this formula we find that the radius of the
6 C nucleus is R ≈ (1·3) (12)
12 1/3 = 3 fm. Similarly,
Nuclear Charge
The model of the nuclear charge distribution
which is the best fit to the experimental scattering
observations for nuclei over a wide range of mass
number is known as the Fermi model and is
represented by the equation
ρ1
ρ(r) = ,
1 + exp. {k (r – c)}
where ρ1 and k are constant and c is the distance
from the centre of nucleus to the point where ρ(r)
has fallen to half its central values. Theoretical
curves based on this model are in good agreement
with this experimental result. The value of c is
found to vary with nucleon number is c = 1·07
A1/3 fm.
r
Charge density

107
the radius of the 47Ag nucleus is 6·2 fm and that
of 92U238 nucleus is 8·1 fm.
The nuclear radius may be estimated from the c
scattering of neutrons and electrons by the
nucleus, or by analysing the effect of the finite Distance from centre
size of the nucleus on nucleus and atomic binding
Fig. : Fermi model for nuclear charge distribution
energies.
Parity—Parity is the fundamental nuclear
Nuclear Density property. In quantum mechanics, the wave
Nucleus is a highly compressed state of mass. function ψ associated with each particle depends
The nuclear density ρN can be calculated from upon the space coordinates (x, y, z). This function
Nuclear mass is not scaler. The parity of the particle decides the
ρN = behaviour of wave function ψ when the sign of x,
Nuclear volume
mN . A mN y, z are changed. If the change of sign of x, y, z
= = does not change ψ, the particle is said to be in
4 4
πγ 30 . A πγ 30 even state or have even parity. If ψ also changes
3 3 with change in sign, the particle is said to be in
1·67 × 10–27 odd state and said to have negative or odd parity.
=
4
× π × (1·3 × 10–15 )3 Nuclear parity is the product of the particle of
3 nuclear constituents. Every nuclear state is
= 1·816 × 1017 kg m–3 characterised by a definite parity. For the same
Certain stars called neutron stars have density nucleus, this may be different for different states.
of this order. The parity of a given nucleus can be given in
710 | Physical Sciences
terms of orbital angular momentum L. In general,
the parity of system of n particles is given by
ΣLn .The parity is even if the sum is positive and
odd if the sum is negative. It can be given as (–1)L.
Nuclear Spin
In addition to the spin angular momentum,
the nucleons have an orbital angular momentum.
The resultant angular momentum of the nucleus
is obtained by adding the spin and orbital
momenta of all the nucleons. This is called
nuclear spin.
LN = 
√ IN (IN + 1)h
Electric Quadrupole Moment
Besides a magnetic moment, a nucleus may
have an electric quadrupole moment. Most nuclei
do possess approximately ellipsoid shape. And the
deviation from spherical body symmetry is
expressed in electric quadrupole moment. It is
defined as
Q = ( )∫
1
e
(3z2 – r2 )ρ dλ,
where ρ is the charge density of nucleus. Q is
actually a measure of the eccentricity of the
ellipsoidal nuclear surface. Evidently Q = 0 for
spherically symmetric charge distribution. A
charge distribution stretched in the –ve direction
(prolate) will give a +ve quadrupole moment
and an oblate distribution will give a –ve
quadrupole moment (see figure). Since, the
expression is divided by the electronic charge, the
dimension of the quadrupole moment is that of an
area. In nuclear physics, area is measured in barn
(1 barn = 10–28 m2).
a b b stability against disintegration. If the binding
a a a energy is large the nucleus is stable. A nucleus
Q=0 Q – ve Q + ve
Statistics
Both the nucleons, i.e., protons and neutrons
have spin 1/2. Nucleons are fermions. If the
nucleon number A of a nucleus is even then
nuclear will follow Bose-Einstein statistic. And in
case of A being odd, the concerned nucleus will
follow Fermi-Dirac statistics.
Natural Radioactivity
Certain elements called radioactive emit
spontaneously invisible radiations which pene-
trate through opaque substances and affects
photographic plates. This phenomenon is called
radioactivity. There are five fundamental radio-
active decays—
(a) α-decay
(b) β-decay
(c) γ-decay
(d) Electron capture
(e) Positron emission.
Artificial Radioactivity
We can convert a stable element into a
radioactive isotope by an artificial disintegration.
In this process, the daughter nuclei have much
shorter half-lives.
13Al + 2 He → 15 p + 0 n
27 4 30 1
e.g.,
(α-particle) Neutron
Binding Energy
The mass of a nucleus is always less than the
sum of masses of its constituents. The decrease in
the mass is known as mass defect (∆m ). The
theoretical explanation for the mass defect is
based on Einstein’s equation E = mc2. When the Z
protons and N neutrons combine to make a
nucleus, some of the (∆m) disappears because it is
converted into its equivalent amount of energy
∆E = (∆m )c2. This energy is called the binding
energy of the nucleus. To disrupt a stable nucleus
into its constituents protons and neutrons, the
minimum energy required is binding energy or
the difference between the energy of the
constituent particles and that of whole nucleus is
defined as the binding energy of the nucleus. The
magnitude of binding energy of a nucleus is its
having the least possible energy, equal to B.E. is
said to be is Ground state. If the nucleus has an
energy E > Emin, it is said to be in excited state.
The case E = 0 corresponds to dissociation of the
nucleus into its constituents nucleus.
If M is the experimentally determined mass of
a nuclide having Z protons and N neutrons,
B.E. = {(Zmp + Nmn) – M}e2.
If B.E. > 0, the nucleus is stable and energy
must be supplied from outside to disrupt it into its
 Physical Sciences | 711

Binding energy per nucleon

constituents. If B.E. < 0, the nucleus is unstable
and it will disintegrate by itself.
Let us see the example of deuteron, which is a

(in MeV)
bound state of a proton and a neutron.
Mass of proton = 1·007276 u
Mass of neutron = 1·008665 u
∴ Mass of proton + Mass of neutron in free
state = 2·01594 u
Mass of deuteron nucleus = 2·013553 u 0 50 100 150 200 250
Mass Number (A)
∴ Mass defect (∆m) = 0·002388 u
B.E. = 0·002388 × 931 = 2·33 MeV Packing Fraction
(Q 1 u = 931 MeV) The ratio between the mass defect (∆m) and
the mass number (A) is called the packing fraction
Thus, when a deuteron is formed from a free (f) :
proton and a free neutron , 2·33 MeV energy is (∆m)
liberated. f = .
A
Proton Since, atomic masses are measured relative to
C-12, the packing fraction for this isotope is zero,
2.33 MeV
packing fraction is a measure of the comparative
stability of the atom.
γ Photon
Neutron Packing fraction
Isotopic mass – Mass number
Conversely, 2·33 MeV must be supplied from = × 104
Mass number
an external source to break a deuteron up into a
proton and a neutron (see figure). This is packing fraction may have a negative or positive
confirmed by γ-ray experiments, where a γ-photon value.
of 2·33 MeV splits the deuteron up into a free For negative packing fraction some mass get
proton and a free neutron. transformed into energy in the formation of that
nucleus, in accordance with Einstein’s equation
Stability of Nucleus and Binding E = (∆ m )c2. Such nuclei are, therefore, more
Energy stable a positive packing fraction would imply a
tendency towards instability. But it is not quite
The binding energy per nucleon is given by correct especially for elements of low atomic
B.E. per nucleon masses.
Total B.E. of a nucleus A plot of packing fraction Vs corresponding
= masses is shown below. It is seen that He, C and
The number of nucleons it contains
O of mass number 4, 12 and 16 respectively, do
The binding energy per nucleon Vs mass not fall on this curve. Their packing fractions are
number curve is shown below. The curve rises smaller and hence these elements are stable. The
steeply at first and then more gradually until it packing fraction of transition elements (around A
reaches a maximum of 8·79 MeV at A = 56 (Iron = 45) have lowest packing fraction i.e., highest
= 26Fe56), then the curve drops slowly to about 7·6
MeV at the highest mass numbers. Evidently, 20
nuclei of intermediate mass are most stable. This
fact suggest that a large amount of energy will be
0
liberated if heavier nuclei can somehow be split
into lighter ones or if light nuclei can somehow be
pointed to form heavier one. The former process is –20
known as Nuclear Fission and later is Nuclear
Fusion. 0 40 80 120 160 200
712 | Physical Sciences

stable and beyond A = 200 the packing fraction helium. The difference in the mass of daughter
becomes positive and increases with increase in nuclei with the parent nuclei is converted into its
mass number increasing the instability of elements. equivalent energy according to Einstein’s Mass-
Elements beyond A = 230 are radio-active and Energy relation
undergo disintegration spontaneously. E = mc 2 .
Nuclear Fission
It is a process in which a heavy nucleus, after
capturing a neutron splits up into two lighter
nuclei of comparable masses, e.g.,
1
0
n
235
92 U → 236
92 U
236
92 U → 144
56 Ba + 89
36Kr + 3 10 n
All fragments undergoes several β-decays
until they reach some stable end product.
A fission yield curve can show the mass
distribution of fission fragments.
Fig. : Detailed process of a proton-proton cycle :
Formation of helium nucleus from four
hydrogen nuclei.

Chain Reaction
If one or more of the emitted neutrons in the
fission process are used to fission other nuclei,
further neutrons are produced. The reaction thus
becomes self sustained and is known as chain
reaction.
The neutron multiplication factor ‘k’ decides
the critically of a nuclear reaction.
No. of neutron in any one generation
k=
No. of neutron in the preceding generation
The fission chain reaction will be critical
when k = 1, subcritical for k < 1 and super critical
for k > 1.
Nuclear Fusion
In this process two or more nuclei combine
together to form a single heavy nucleus. e.g., four
hydrogen nuclei combine to form a nucleus of
1H
1 + 1 H1 + 1 H1 + 1 H1 → 2 He 4 + Energy
This fusion is the prime source of steller
energy. A large amount of energy is released by
nuclear fusion. But unlike nuclear fission one have
to overcome the Coulombian potential (repulsive)
barrier between the lighter nuclei to fuse them.
This condition named as Thermonuclear reactions
is not feasible on earth under normal conditions.
Controlled Thermonuclear Reactions
A large amount of energy is released in a
thermonuclear process in a fraction of second. If
this energy can be controlled, we can use it for
constructive purposes. Also we need very high
temperature and great pressure to maintain the
fusion reaction. At this high temperature the gas is
highly ionised and is called plasma. Plasma cannot
be confined in any mechanical container because
it will melt. Scientists are designing a magnetic
bottle in which the plasma can be confined with
the help of variable magnetic field.

NATURE OF NUCLEAR FORCES
Interaction between Energetic Particles
and Matter
(a) Heavy charged particles—A heavy
charged particle (like a proton, α-particle or
fission fragments) has a fairly definite range in a
gas, liquid or solid. The particle losses energy
primarily by the excitation and ionization of atoms
in its path. The energy loss occurs in a large
number of small increments. The primary particle
has such a large momentum that its direction is
usually not seriously changed during the slowing
process. Eventually after the loss of its total
energy, it comes to rest.The distance traversed is
called the range of the particle.
The energy loss per unit length
–dE
dx ( ) is
–dE
called the stopping power. The gradient or the
dx
rate at which the heavy particle of charge Ze and
speed v loses energy in an absorber of atomic
number Z which contains N atoms per unit
volume whose average ionization energy is I, is
given by
–dE Z2 e4NZ
= known as linear absorption coefficient. The mass
dx 4πε20 mv20
  µ
( 2m0 v2
) (1 – vc ) – vc 
2 2
 log – log
I 2 2 ρ
where m0 is rest mass of an electron.
The range can be calculated by integrating
equation (1) over the energies of the particle
∫
T
R =
0 (–dE
dx )
–1
dE
(b) Electrons—Electrons interact through
Coulomb scattering from atomic electrons, just
like heavy charged particles. There are however, a
number of important differences—
Unit–11
(i) Electrons travel at relativistic speeds.
(ii) Electron follow erractic path due to large
deflection in collision with other electrons. The
range will, therefore, be very different from the
length of the path that the electron follow.
(iii) Very energetic electrons (E > 1 MeV)
lose an appreciable fraction of their energies by
the production of continuous X-rays. (also called
Bremsstrahlung). The cross-section for this process
increases with increasing E.
(c) Absorption of γ -rays—The interaction of
γ-rays with matter is completely different from
that of charged particles such as α or β particles.
γ-rays are extremely penetrating so they are able
to pass through considerable thickness of matter.
If radiation of intensity I is incident upon an
absorbing layer of thickness dx, the amount of
radiation absorbed dI is proportional to dx and I.
Hence,
dI = –µIdx
or I = I0 e – µ x
Here µ is the constant of proportionality
which is a characteristic property of medium,
absorption coefficient µm may be obtained by
…(1) dividing µ by the density of medium µm = .
At low energies three processes takes place.
Photoelectric effect
hν e–
…(2)
hν e–
hν
Compton effect
e+–
e Pair production
hν
Before interaction After interaction
756 | Physical Sciences

(i) Photoelectric effect—In photoelectric energy, then Compton effect dominates. Pair pro-
effect, the γ-rays knock-out electrons from inside duction becomes increasingly likely as the more
the atom of absorbing material. This results in the the photon energy exceeds the threshold of 1·02
ionisation of the atoms and the emission of MeV. A graph of linear attenuation coefficient for
fluorescent radiation. Einstein’s equation for these photon as a function of photon energy is shown
photoelectron will be below in which the contribution to µ is the above
1 three effects is shown.
m v2K = hν – hνK = hν – WK
2 Types of Nuclear Reactions
1 Following are some of the main types of
m v2L = hν – hνL = hν – WL
2 nuclear reactions :
………………………………...
………………………………... (i) Elastic scattering—In this case, the inci-
dent particle strikes the target nucleus and leaves
where hν is the photon energy, νK , νL represent without loss of energy, but its direction may
the velocities of the photoelectrons arising in K, L change, e.g., scattering of an α-particle by a gold
shells and WK, WL …… are the binding energies foil.
of K, L …… shells. 197 + He 4 → Au197 + He 4
79Au 2 79 2
(ii) Pair production—In pair production, the The target nucleus remains unaffected.
photon disappears and gets converted into a pair
of electron positron pair. This process can take (ii) Inelastic scattering—In this case, the
place only when the energy of incident photon is incident particle loses a part of its energy in
greater than 2m 0 c2. The pair production cannot exciting the target nucleus to a higher allowed
occur in free space and usually takes place is the energy level. The excited nucleus later decays to
presence of nuclear field. Nucleus recoils in the the ground state, radiating the excess energy in the
process to conserve momentum. But the K.E. form of a γ-ray photon e.g.,
3 Li + 1 H → 3 Li + 1 H
7 1 7* 1
carried away by the nucleus is negligible. Any
energy more than 2m0c2 is shared as K.E. by the 3 Li → 3 Li + γ
7* 7

product particles. (iii) Radioactive capture—Here the incident

(iii) Compton effect—It is an elastic scatter- particle is captured by the target nucleus and a
ing process in which the photon imparts energy to new nucleus is formed. The new nucleus, in
an electron. When a photon of energy hν strikes general, has a considerable excess of energy and
the perfectly free electron (at rest), the photon decays with the emission of one or more γ-ray
with reduced energy hν′ is scattered at an angle θ photons, example
with the direction of incident photon and the 6C + 1H → 7N
12 1 13*
electron recoils at an angle φ. The energy absor- 13* → N13 + γ
7N 7
bed by these Compton electrons is only a small
fraction of total energy of the incident γ-rays. (iv) Disintegration—In disintegration reac-
tions, the incident particle is absorbed by the tar-
Linear attenuation coefficient (cm–1)

1.8 - get nucleus and the ejected particle is a different

1.6 - one. The composition of the resultant nucleus is
1.4 -
1.2 -
also different from the parent nucleus. An exam-
1.0 -
Total
otal
ple is the disintegration of beryllium by α-particle
n
.8 - uctio producing neutrons
Prod
.6 - Pair
4 Be + 2 He → 6 C + 0 n
.4 - 9 4 12 1
.2 - CE effect Compton Scattering
scattering
0 (v) Photodisintegration—When target nuc-
5 10 15 20 25
leus are bombarded with radiations, the resulting
compound nuclei are usually formed in excited
Photon energy (MeV) states. These nuclei generally get rid of the excess
excitation energy through neutrons emission, e.g.,
At low energy, the photoelectric effect is the
1H + γ → 1H + 0n
2 1 1
chief mechanism of energy loss. The importance
of photoelectric effect decreases with increase in This requires a photon of 2·225 MeV.
 Physical Sciences | 757

Conservation Laws Energy Balance in Nuclear Reaction

The conservation laws may be explained by
reference to some specific nuclear reaction, say
5 B + 2 He → (7 N ) → 1H + 6 C + Q
10 9 14 1 13
(i) Conservation of charge—Total charge is
conserved in energy type of nuclear reaction. In
B 10 (α, p)C13 there are seven protons initially, also
seven in the compound nucleus and in the
products of the reaction.
(ii) Conservation of nucleons—The total
number of nucleons entering and leaving the
reaction is constant. In B10(α, p)C13, we find 14
nucleons at each stage of the reactions.
(iii) Conservation of mass energy—In
nuclear reaction neither kinetic energy nor rest
mass is conserved by itself. But the mass and
energy in combined form are conserved.
(iv) Conservation of parity—The net parity
before the reaction must equal the net parity after
the reaction.
(v) Linear momentum, angular momentum,
spin and isotopic spin are the other quantities
which are also conserved in a nuclear reaction.
Quantity not Conserved
The prominent physical quantities which are
not conserved in nuclear reactions are the magnetic
dipole moments and the electric quadrupole
moments of the reacting nuclei. These moments
depend upon the internal distribution of mass,
charge and current within the nuclei involved and
are not subject to conservation law.
Elementary Particles
and the Q-Value
In all nuclear reactions, the total sum of
mass and energy is conserved. Thus taking the
equation A + B → X + Y the target is supposed to
be at rest and let its mass be m1. The projectile has
a mass m 2 and K.E. = E 2 .
The product nucleus has a mass m 3 and K.E.
E3 . The outgoing particle has a mass m4 and K.E.
E4 . The equation representing the conservation of
total-energy is written as :
m1c2 + m2c2 + E2 = m3c2 + E3 + m4c2 + E4
Let Q = E3 + E4 – E2 = (m1 + m2 – m3 – m4)c2 .
The quantity is called the energy balance or
the Q-value of the reaction. If Q is positive, the
reaction is exothermic (or exoergic). Now the K.E.
of products of transmutation is greater than the
K.E. of the reactants and energy is released in the
process. If Q is negative, the reaction is
endothermic. The equation for a nuclear reaction
is written as
A+B→X+Y
where Q may be positive or negative.
Elementary Particles
Elementary particles are the quanta of
respective fields. In another definition, elementary
particles cannot be further subdivided. Following
is a table to classify the elementary particles.

Bosons Fermions
(Integral spin particles) (Half integral spin particles)

Massless Mesons Leptons Baryons

Photon Graviton e e µ µ ν ν νµ νµ

Leptons Kaons Pions Nucleon Hyprons

k k k k π π π p p n n

Lambda Sigma Cascade Omega

hyperon hyperon hyperon hyperon

Σ Σ Σ
758 | Physical Sciences

Fundamental Interactions (Particle 4. Weak interactions—These interactions

Interactions) are responsible for the decay of strange and non-
strange particles. Their strength is weak but better
1. Gravitational interactions—These are than gravitational interaction. The quantum of
always attractive inverse square law forces. These
these interaction is charged vector boson (W and
are independent of size, charge, velocity or spin Z) and characteristic time is 10–10 sec.
but its depends on inertia. Its field particle is
graviton (Not detected experimentally so far) and Attempts are being made to combine these
characteristic time 10–16 sec. These are the most interactions and develope a successful theory that
feeble forces of nature. unites all the ways in which particles of matter
2. Electromagnetic interaction—These are interact with each other.
charge dependent inverse square law forces. These
may be attractive or repulsive in nature depending Electricity Magnetism
upon the polarity of charges. These depend on the
velocity and magnitude of charge. EM
Electro- Weak
interactions are around 1035 times stronger to magnetic interaction
gravitational interaction. The quantum EM field is
photon and characteristic time is 10–20 sec.
Terrestrial Astronomical Electro Strong
3. Strong or nuclear interactions—These gravity gravity Weak Interaction Interaction
are the strongest forces of nature ever known.
These are short range charge independent forces. Gravitational Grand Unified
These forces cannot be described in terms of Interaction Interaction
strength distance relation easily. These are spin
Task yet to be
dependend forces. The quanta of field are pions Universal completed
and kaons and characteristic time 10–23 sec. Interaction

Properties of three Generations of Quarks and Leptons

Generation Name of Quark Symbol Charge e Strangeness Charm
2
1. Up u + 0 0
3
1
Down d – 0 0
3
2
2. Charm c + 0 +1
3
1
Strange s – –1 0
3
2
3. Top t + 0 0
3
Bottom (Beauty) 1
b – 0 0
3

Generation Name of Lepton Symbol Charge C

1. Electron e– –1
e-nutrino νe 0
2. Muon µ– –1
µ-neutrino νµ
0

3. Tau τ– –1
τ-neutrino ντ
0
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