𝟐𝑿 : 2𝑋 is used in functional analysis, Measure Theory etc to denote power set 𝑃(𝑋) of all

subsets of 𝑋.
α-CUT OF A FUZZY SET: An α-cut of a fuzzy set 𝐴 is a crisp set 𝐴, that contains all the
elements in 𝑈 that have membership values in 𝐴 greater than or equal to a, that is,
𝐴𝛼 = {𝒂 ∇ 𝑈; 𝝁𝑨 (𝑥) ≥ 𝛼}
ς-ALGEBRA: Let X be a set. A ς-algebra on 𝑋 is a collection of subsets of 𝑋, say 𝑅 ⊆ 2𝑋 ,
such that

1. 𝑋 ∇ 𝑅;

2. if 𝐴, 𝐵 ∇ 𝑅, then 𝐴 ∖ 𝐵 ∇ 𝑅;

3. if (𝐴𝑛 ) is any sequence in 𝑅, then ⋃𝑛 𝐴𝑛 ∇ 𝑅.

If we drop the first condition from the definition we got a ς-ring. Note, that in the third
condition we admit any countable unions. The empty set belongs to any ς-algebra. For
a 𝜍-algebra 𝑅 and 𝐴, 𝐵 ∇ 𝑅, we have 𝐴 ∩ 𝐵 = 𝑋\ 𝑋\(𝐴 ∩ 𝐵) = X ((X\A)⋃(X\B) ) ∇ R .
Note that

 As the intersection of a family of 𝜍 −algebras is again a 𝜍 −algebra.

 The power set 2𝑋 is a 𝜍 −algebra.
 Given any collection ⊆ 2𝑋 , there is a 𝜍 −algebra 𝑅 such that 𝐷 ⊆ 𝑅, such that if 𝑆
is any other 𝜍 −algebra, with 𝐷 ⊆ 𝑆, then 𝑅 ⊆ 𝑆. We call R the ς-algebra
generated by 𝐷.

ς-LOCALLY FINITE BASE: A ς-locally finite base is a base which is a union of countably
many locally finite collections of open sets.

Ɛ-Net: A sub set S of a metric space X is called ɛ-net if for any point in X there is a point
in S on the distance ≤ ɛ. This is different from topological nets which oversimplify
limits.

A
A-SET: It is the continuous image of a Borel set. Since any Borel set is a continuous
image of the set of irrational numbers, an A-set can be defined as a continuous image of
the set of irrational numbers.
 A countable intersection and a countable union of A-sets is an A-set.
 Any A-set is Lebesgue-measurable.
 The property of being an A-set is invariant relative to Borel-measurable
mappings, and to A-operations. Moreover, for a set to be an A-set it is necessary
and sufficient that it can be represented as the result of an A-operation applied to
a family of closed sets.
 Any uncountable A-set topologically contains a perfect Cantor set.
ABACUS: Abacus is a counting device consisting of rods on which beads can be moved
so as to represent numbers.
abc CONJECTURE: (1) The abc conjecture of Masser and Oesterlé efforts to state as much
as possible about repeated prime factors in an equation a + b = c. For example 3 + 125
= 128 but the prime powers here are special.
(2) 𝑎𝑚 + 𝑏 𝑛 = 𝑐 𝑝 has no solution in integers if a, b, c >2
ABEL DIFFERENTIAL EQUATION: The ordinary differential equation

𝑦 ′ = 𝑓0 𝑥 + 𝑓1 𝑥 𝑦 + 𝑓2 𝑥 𝑦 2 + 𝑓3 𝑥 𝑦 3 (Abel's DE of the first kind) or

(𝑔0 (𝑥) + 𝑔1 (𝑥)𝑦)𝑦′ = 𝑓0 𝑥 + 𝑓1 𝑥 𝑦 + 𝑓2 𝑥 𝑦 2 + 𝑓3 𝑥 𝑦 3 (Abel's DE of the second

kind). Abel's differential equations of the first kind represent a natural generalization of
the Riccati equation.

If 𝑓1 ∇ 𝐶(𝑎, 𝑏) and 𝑓2 , 𝑓3 ∇ 𝐶1 (𝑎, 𝑏) and 𝑓3 (𝑥) ≠ 0 for 𝑥 ∇ [𝑎, 𝑏], then Abel's differential
equation of the first kind can be reduced to the normal form 𝑑𝑧/𝑑𝑡 = 𝑧 3 + 𝛷(𝑡) by
substitution of variables. In the general case, Abel's differential equation of the first kind
cannot be integrated in closed form, though this is possible in special cases.
If 𝑔0 , 𝑔1 ∇ 𝐶1 (𝑎, 𝑏) and 𝑔1 (𝑥) ≠ 0, 𝑔0 (𝑥) + 𝑔1 (𝑥)𝑦 ≠ 0, Abel's differential equation of
the second kind can be reduced to Abel's differential equation of the first kind by
substituting
𝑔0 (𝑥) + 𝑔1 (𝑥)𝑦 = 1/𝑧.
ABELIAN CATEGORY: An Abelian category is an abstract mathematical category which
displays some of the characteristic properties of the category of all abelian groups.
ABELIAN DIFFERENTIALS: An Abelian differential on a closed Riemann surface ℛ is, by
definition, a complex differential form 𝑤 = 𝑎(𝑧)𝑑𝑧, where 𝑎(𝑧) is a meromorphic
function of a local parameter 𝑧. Such a differential is said to be of the first kind if 𝑎(𝑧) is
holomorphic, of the second kind if the residue vanishes everywhere, and of the third
kind otherwise.
ABELIAN FUNCTION: An inverse function of an abelian integral. Abelian functions have
two variables and four periods. They are a generalization of elliptic functions, and, are
also called hyper elliptic functions.

ABELIAN GROUP: A group G with binary operation ∗ is said to be an abelian group if

𝑎 ∗ 𝑏 = 𝑏 ∗ 𝑎 ∀ 𝑎, 𝑏 ∇ 𝐺 i.e. G is commutative under the operation ∗. For example, the set
of real numbers and the set of all rational numbers form an Abelian group with respect
to usual addition, and the set of nonzero real numbers and the set of all nonzero rational
numbers form an Abelian group with respect to usual multiplication. The theory of
Abelian groups has its origins in number theory, and is now extensively applied in many
modern mathematical theories. Thus, the duality theory of characters for finite Abelian
groups has been considerably extended to the duality theory of locally compact Abelian
groups. The development of homological algebra has made it possible to solve a whole
series of problems in Abelian groups, such as classifying the set of all extensions of one
group by another. The theory of modules is closely connected with Abelian groups
regarded as modules over the ring of integers. Many results in the theory of Abelian
groups can be applied to the case of modules over a principal ideal ring. Owing to their
relative simplicity and to the fact that they have been very thoroughly studied (which is
confirmed, for instance, by the solvability of the elementary theory of Abelian groups),
and to the availability of a sufficient variety of objects, Abelian groups serve as a
constant source of examples in various fields of mathematics.

Examples of Abelian groups:

 All cyclic groups are Abelian, in particular, the additive group of integers.
 All direct sums of cyclic groups are Abelian. Also the additive group of rational
numbers Q is Abelian; it is moreover a locally cyclic group, i.e. a group in which
all finitely generated subgroups are cyclic.
 The free composition in the variety of Abelian groups coincides with the direct
sum. A free Abelian group is a direct sum of infinite cyclic groups.
  Every subgroup of a free Abelian group is free Abelian.
 The set of all elements of finite order in an Abelian group forms a subgroup,
which is called the torsion subgroup of the Abelian group.
 The quotient group of an Abelian group by its torsion subgroup is a group
without torsion.
 Every Abelian group is an extension of a torsion-free group by a torsion group.
This extension does not always split, i.e. the torsion group is usually not a direct
summand.
 An Abelian torsion group in which the order of every element is a power of a
fixed prime number p is said to be primary with respect to p (the term p-group is
used in general group theory).
 Every torsion group splits uniquely into a direct sum of primary groups that
correspond to distinct prime numbers.

A complete classification is known for finitely generated Abelian groups. This is given by
the fundamental theorem of finitely generated Abelian groups:

Every finitely generated Abelian group is a direct sum of finitely many non-split cyclic
subgroups some of which are finite and primary, while the others are infinite. In
particular, finite Abelian groups split into a direct sum of primary cyclic groups. Such
splittings are, in general, not unique, but any two splitting of a finitely generated Abelian
group into direct sums of non-split cyclic groups are isomorphic, so that the number of
infinite cyclic summands and the collection of the orders of the primary cyclic
summands do not depend on the splitting chosen. These numbers are called invariants
of the finitely generated Abelian group.

 Two finitely generated Abelian groups are isomorphic if and only if they have the
same invariants.

 Each subgroup of a finitely generated Abelian group is itself finitely generated.

 Not every Abelian group is a direct sum of (even infinitely many) cyclic groups.
 Any subgroup of an Abelian group that is a direct sum of cyclic groups itself is
such a direct sum.
A finite set of elements 𝑔1 , … , 𝑔𝑘 in an Abelian group is called linearly dependent if there
𝑘
exist integers 𝑛1 , … , 𝑛𝑘 , not all equal to zero, such that 𝑖=1 𝑛𝑖 𝑔𝑖 = 0. If such numbers do
not exist, the set is said to be linearly independent. An arbitrary collection of elements
of A is said to be linearly dependent if there exists a finite linearly dependent sub
collection. An Abelian group that is not a torsion group has maximal linearly
independent sets. The cardinality of all maximal linearly independent collections of
elements is the same and is called the rank of the given Abelian group. The rank of a
torsion group is considered to be zero. The rank of a free Abelian group coincides with
the cardinality of a set of free generators of it.
ABEL’S IMPOSSIBILITY THEOREM: In general, polynomial equations higher than
fourth degree are incapable of algebraic solution in terms of a finite number of
additions, multiplications, and root extractions.

ABEL’S INEQUALITY: If ⌌𝑎𝑛 ⌍ and ⌌𝑏𝑛 ⌍ are sequences such that for some 𝑚 ∇ 𝑁, and
𝑛
𝛼 ∇ 𝑅 +, 𝑏1 ≥ 𝑏2 ≥ 𝑏3 ≥ . . . ≥ 𝑏𝑚 −1 ≥ 𝑏𝑚 ≥ 0 and |𝑠𝑛 | ≅ 𝑖=1 𝑎𝑖 ≤𝛼 for
𝑛
𝑛 = 1, 2, . . . , 𝑚, then 𝑖=1 𝑎𝑖 𝑏𝑖 ≤ 𝛼𝑏1 .
ABEL’S IRREDUCIBILITY THEOREM: If one root of the equation 𝑓(𝑥) = 0, which is
irreducible over a field K, is also a root of the equation 𝐹(𝑥) = 0 in K, then all the roots
of the irreducible equation 𝑓(𝑥) = 0 are roots of 𝐹(𝑥) = 0. Equivalently, 𝐹(𝑥) can be
divided by 𝑓(𝑥) without a remainder, 𝐹(𝑥) = 𝑓(𝑥)𝑔(𝑥), where 𝑔(𝑥) is also a polynomial
over K.

ABEL, NIELS HENRIK (1802–29): Norwegian mathematician who, at the age of 19,
proved that the general equation of degree greater than 4 cannot be solved
algebraically. He was also responsible for fundamental developments in the theory of
algebraic functions. He died in some poverty at the age of 26, just a few days before he
would have received a letter announcing his appointment to a professorship in Berlin.

ABEL'S CONTINUITY THEOREM: If the power series

∞
𝑘
𝑆 𝑧 = 𝑎𝑘 𝑧 − 𝑏
𝑘=0

converges at a point 𝑧0 on the boundary of the disc of convergence, then it is a

continuous function in any closed triangle with vertices 𝑧0 , 𝑧1, 𝑧2 , where 𝑧1, 𝑧2 are
located inside the disc of convergence. In particular

lim 𝑠 𝑧 = 𝑠(𝑧0 )
𝑧→𝑧0
This limit always exists along the radius: The series
∞
𝑘
𝑆 𝑧 = 𝑎𝑘 𝑧 − 𝑏
𝑘=0

converges uniformly along any radius of the disc of convergence joining the
points 𝑏 and 𝑧0 . This theorem is used, in particular, to calculate the sum of a power
series which converges at the boundary points of the disc of convergence.

ABEL’S METHOD OF SUMMATION: If the radius of convergence of the power series

∞ 𝑛 ∞ 𝑛
𝑛=1 𝑢𝑛 𝑟 is ≥ 1 and 𝑛=0 𝑢𝑛 𝑟 → 𝑠 𝑎s 𝑟 → 1, then 𝑢𝑛 is said to be summable by
Abel’s method or A-summable) to 𝑠, and we write 𝑢𝑛 = 𝑠(𝐴). The transformation
matrix is denoted by 𝐴, and the transformation is called Abel’s method of summation.

ABEL’S PARTIAL SUMMATION: Let 𝑎0 , 𝑎1 , 𝑎2 , … 𝑎𝑛𝑑 𝑏0 , 𝑏1 , 𝑏2 , … be arbitrary

sequences, and put 𝐴𝑛 = 𝑎0 + 𝑎1 + 𝑎2 + ⋯ + 𝑎𝑛 for 𝑛 ≥ 0. Then the following formula
of Abel’s partial summation holds;

n+k n+k

a v bv = Av bv − bv+1 − An bn−1 +an+k bn+k−1

v=n+1 v=n+1

For any n ≥ 0 and any k ≥ 1; this formula also holds for n = −1 if we put 𝐴−1 = 0.

Able’s partial summation enables us to deduce a number of tests of convergence for

series of the form 𝑎𝑛 𝑏𝑛 . In particular, the following criteria are easy to apply:

(1) 𝑎𝑛 𝑏𝑛 is convergent if 𝑎𝑛 is convergent and if the sequence 𝑏𝑛 is montone

and bounded Abel’s test .
(2) 𝑎𝑛 𝑏𝑛 is convergent if the sequence 𝑠𝑛 of partial sums of 𝑎𝑛 is bounded and
if 𝑏𝑛 is monotone and converges to zero Dirichlet’s test .
(3) 𝑎𝑛 𝑏𝑛 is convergent if (𝑏𝑛 − (𝑏𝑛+1 ) is absolutely convergent and if 𝑎𝑛 is (at
least conditionally) convergent (test of du Bois- Reymond and Dedekind).

ABEL’S TEST: Abel’s test is a test for the convergence of an infinite series which states
that if 𝛴𝑎𝑛 is a convergent series, and {𝑏𝑛 } is monotonically decreasing, then 𝛴𝑎𝑛 𝑏𝑛 is
also convergent.
ABEL’S THEOREM (ALGEBRAIC EQUATIONS): Formulas expressing the solution of an
arbitrary equation of degree in terms of its coefficients using radicals do not exist for
any 𝑛 ≥ 5 . The theorem was proved by N.H. Abel in 1824. Abel's theorem may also be
obtained as a corollary of Galois Theory, from which a more general theorem follows:
For any 𝑛 ≥ 5 there exist algebraic equations with integer coefficients whose roots
cannot be expressed in terms of radicals of rational numbers.

ABEL'S THEOREM ON POWER SERIES: If the power series

∞
𝑘
𝑆 𝑧 = 𝑎𝑘 𝑧 − 𝑏 (∗)
𝑘=0

where 𝑎𝑘 , 𝑏, 𝑧 are complex numbers, converges at 𝑧 = 𝑧0 , then it converges absolutely

and uniformly within any disc 𝑧 − 𝑏 ≤ 𝜌 of radius 𝜌 ≤ 𝑧0 − 𝑏 and with centre at 𝑏.
The theorem was established by N.H. Abel. It follows from the theorem that there exists
a number 𝑅 ∇ [0, ∞) such that if 𝑧 − 𝑏 ≤ 𝑅 the series is convergent, while
if 𝑧 − 𝑏 ≤ 𝑅 the series is divergent. The number 𝑅 is called the radius of convergence
of the series (*), while the disc 𝑧 − 𝑏 ≤ 𝑅 is known as the disc of convergence of the
series (*).

ABOVE (GREATER THAN): The limit of a function at 𝑎 from above is the limit of 𝑓(𝑥) as
𝑥 → 𝑎 for values of 𝑥 > 𝑎. It is of particular importance when 𝑓(𝑥) has a discontinuity
at 𝑎, i.e. where the limits from above and from below do not coincide. It can be written
as 𝑓(𝑎+) or lim𝑥→𝑎+ 𝑓(𝑥).
ABRIDGED MULTIPLICATION: A method of multiplication where digits may be dropped
in each step of the calculation, in order to simplify the calculations, while maintaining
the desired level of accuracy, usually expressed in terms of decimal places.

ABSCISSA: Abscissa in two dimensional plane means x-coordinate. The abscissa of the
point (a, b) in Cartesian coordinates is a.
ABSOLUTE CONTINUITY OF LEBESGUE INTEGRAL: Let 𝑓 ∇ 𝐿1 (𝑋). Then for any 𝜀 > 0
there is a 𝛿 > 0 such that

∫ 𝑓 𝑑µ < 𝜀 𝑖𝑓 µ 𝐴 < 𝛿.
𝐴

Let 𝑋 and 𝑌 be spaces, and let 𝑆 and 𝑇 be semirings on 𝑋 and 𝑌 respectively and µ and 𝜈
be measures on 𝑆 and 𝑇 respectively. If µ and 𝜈 are 𝜍 −additive, then the product
measure 𝜈 × µ is 𝜍 −additive as well.
Let 𝐶 ∇ 𝐿(𝑋 × 𝑌). For almost every 𝑥 ∇ 𝑋 the set 𝐶𝑥 = {𝑦 ∇ 𝑌: (𝑥, 𝑦) ∇ 𝐶} is
𝜈 −measurable and 𝜈(𝐶𝑥 ) = 𝑓𝐶 (𝑥).

Any measurable set is up to a set of zero measure can be received from elementary sets
by two monotonic limits.

Let µ and 𝜈 are 𝜍 −finite measures and 𝐶 be a µ × 𝜈 measurable set 𝑋 × 𝑌. We define

𝐶𝑥 = {𝑦 ∇ 𝑌: (𝑥, 𝑦) ∇ 𝐶}. Then for µ −almost every 𝑥 ∇ 𝑋 the set 𝐶𝑥 is 𝜈 −measurable,
function 𝑓𝐶 (𝑥) = 𝜈(𝐶𝑥 ) is µ −measurable and (µ × ν)(C) = ∫ fc dμ, where both parts may
have the value +∞.

ABSOLUTE CONVERGENCE: A series is said to be absolutely convergent if it still

converges when all the terms are replaced by the corresponding absolute values. The
series 1 − (1/2)2 + (1/3)3 − (1/4)4 + . ..is an example of absolute convergence. This
is because 1 − (1/2)2 + (1/3)3 − (1/4)4 + . .. is also convergent.

ABSOLUTE ERROR: The difference between the measured value or the approximation of
a quantity and the true value. When 0.7 is used as an approximation for 0.709, the error
equals (0.7-0.709), which is -0.009. In this case, the absolute error is |0.7-0.709|, which
is 0.009.
ABSOLUTE EXTREMUM: It is also known as an absolute maximum or an absolute
minimum.
ABSOLUTE FREQUENCY: Absolute frequency is the number of occurrences of an event.
For example, if a die is rolled 50 times and 8 sixes are observed, the absolute frequency
of sixes is 8 and the relative frequency is 8/50.
ABSOLUTE GEOMETRY: The term Absolute Geometry had been introduced by J. Bolyai
in 1832. Absolute Geometry is derived from the first four of Euclid's postulates. The first
28 Propositions from Elements, I belong to Absolute Geometry. Euclid apparently made
a conscientious effort to see how far he can reach without invoking his Fifth postulate.
All theorems of Absolute Geometry are automatically true in the geometries of Euclid,
Lobachevsky and Riemann since those three only differ in their treatment of the Fifth
postulate.
ABSOLUTE INEQUALITY: It is an inequality that contains absolute value(s). Example: |7x
- 6| < 1.
ABSOLUTELY CONTINUOUS CHARGE: Let 𝑋 be a set with 𝜍 −algebra 𝑅 and 𝜍 −finite
measure µ and finite charge 𝜈 on 𝑅. The charge 𝜈 is absolutely continuous with respect
to µ if µ(𝐴) = 0 for 𝐴 ∇ 𝑅 implies 𝜈(𝐴) = 0. Two charges 𝜈1 and 𝜈2 are equivalent if two
conditions | 𝜈1 |(𝐴) = 0 and | 𝜈2 |(𝐴) = 0 are equivalent.

ABSOLUTELY CONVERGENT SERIES: A series 𝑢𝑛 is said to be absolutely convergent if

the series 𝑢𝑛 is convergent.
ABSOLUTE MAXIMUM: The absolute maximum point for a function 𝑦 = 𝑓 (𝑥) is the
point where 𝑦 has the largest value on an interval 𝐼. If the function 𝑦 = 𝑓 (𝑥) is
differentiable, the absolute maximum of the function 𝑦 = 𝑓 (𝑥) will either be a point
where there is a horizontal tangent so the derivative is zero, or a point at one of the
ends of the interval.
ABSOLUTE MINIMUM: The absolute minimum point for a function 𝑦 = 𝑓 (𝑥) is the point
where 𝑦 has the smallest value on an interval 𝐼. If the function 𝑦 = 𝑓 (𝑥) is
differentiable, then the absolute minimum of the function 𝑦 = 𝑓 (𝑥) will either be a
point where there is a horizontal tangent so the derivative is zero, or a point at one of
the ends of the interval.

ABSOLUTE NUMBER: A real value with its sign ignored - the result of applying
the modulus function to a value.

ABSOLUTE TERM: The constant term in an expression such as a polynomial.

ABSOLUTE VALUE: The absolute value of a real number 𝑎, written as 𝑎 , is:

𝑎; 𝑖𝑓 𝑎 ≥ 0
𝑎 =
−𝑎; 𝑖𝑓 𝑎 ≤ 0
The following properties hold:
(i) |𝑎𝑏| = |𝑎||𝑏|.
(ii) |𝑎 + 𝑏| ≤ |𝑎| + |𝑏|.
(iii) |𝑎 − 𝑏| ≥ ||𝑎| − |𝑏||.
(iv) For 𝑎 > 0, |𝑥| ≤ 𝑎 if and only if −𝑎 ≤ 𝑥 ≤ 𝑎.
Figure given below illustrates the absolute value function in nbd of zero.
Absolute value of a real number is always positive or zero. If we represent all the real
numbers on a number line, we can think of the absolute value of a number 𝑎 as being
the distance from origin (zero) to that number 𝑎. The absolute value of a complex
number 𝑧 = 𝑎 + 𝑖𝑏 is denoted and defined as

𝑧 = 𝑎2 + 𝑏 2 .
It is to be noted that absolute value of a complex number 𝑧 is again the distance from
origin (zero) to that number 𝑧.
In a similar fashion, we can define the absolute value of a number 𝑎 ∇ 𝑅 𝑛 as
𝑎 = 𝑎1 2 + 𝑎2 2 + − − − − +𝑎𝑛 2 , where 𝑎 = 𝑎1 , 𝑎2 , − − −𝑎𝑛 .
It is to be noted that absolute value of a number 𝑎 ∇ 𝑅 𝑛 is again the distance from origin
(zero) to that number 𝑎 in n dimensional space.
ABSORPTION LAWS: For all sets A and B, we have
𝐴 ∩ (𝐴 ∪ 𝐵) = 𝐴 𝑎𝑛𝑑 𝐴 ∪ (𝐴 ∩ 𝐵) = 𝐴.
These are said to be the absorption laws.
ABSTRACT ALGEBRA: It is the area of mathematics concerned with algebraic structures
involving sets of elements with particular operations satisfying certain axioms such as
groups, rings and fields.
ABSTRACTION: Abstraction is the process of making a general statement which
summarizes what can be observed in particular instances. For example, we can say that
𝑥 3 < 𝑥 for 0 < 𝑥 < 1 and 𝑥 3 > 𝑥 for 𝑥 < 0 𝑜𝑟 𝑥 > 1. Mathematical theorems are
essentially abstraction of concepts to a higher level.
ABSTRACT MANIFOLD: An abstract manifold (or just a manifold) of dimension m, is a
Hausdorff topological space 𝑀, equipped with an m-dimensional smooth atlas.
Compatible atlases are regarded as belonging to the same manifold (the precise
definition is thus that a manifold is a Hausdorff topological space equipped with a
smooth structure). A chart on 𝑀 is a chart from any atlas compatible with the structure.
It is often required of an abstract manifold that it should have a countable atlas.
ABSTRACT SUBMANIFOLD: Let 𝑀 be an abstract manifold. An abstract submanifold of 𝑀
is a subset 𝑁 ⊂ 𝑀 which is an abstract manifold on its own such that:
(i) the topology of 𝑁 is induced from 𝑀,
(ii) the inclusion map 𝑖: 𝑁 → 𝑀 is smooth, and
(iii) the differential 𝑑𝑖𝑝 : 𝑇𝑝 𝑁 → 𝑇𝑝 𝑀 is injective for each 𝑝 ∇ 𝑁.
In this case, the manifold 𝑀 is said to be ambient to 𝑁. In particular, since 𝑑𝑖𝑝 is injective,
the dimension of 𝑁 must be smaller than or equal to that of 𝑀.
ABUNDANCE: The abundance of a number 𝑛 is the quantity

𝐴 𝑛 = 𝜍(𝑛) 2𝑛,

where 𝑜(𝑛) is the divisor function. Kravitz has conjectured that no numbers exist whose
abundance is an odd square.

The following table lists special classifications given to a number n based on the value of
A(n).

A(n) Number
<0 deficient number
-1 almost perfect number
0 perfect number
1 Quasi-perfect number
>1 abundant number
ABUNDANT NUMBER: Abundant number is a positive integer that is smaller than the
sum of its positive divisors, not including itself, For example, 12 is divisible by 1, 2, 3, 4
and 6, and 1 + 2 + 3 + 4 + 6 = 16 > 12.
AC: It is the abbreviation for Axiom of Choice.
ACCELERATION: The acceleration of a moving object measures the rate of change in its
velocity with respect to time. For example, if a car increases its velocity from 20 to 45
45−20
meters/ second in 8 seconds, its acceleration would be = 3.125 meters/ second2
8

to read as 3.125 meters/ second-squared.

If 𝑦(𝑡) represents the position of the moving object as a function of time, then the
𝑑𝑦
derivative represents the velocity of the object, and the second derivative
𝑑𝑡
𝑑2𝑦
represents the acceleration.
𝑑𝑡 2

Newton obtained that, if F represents the force acting on a moving object and m
represents its mass, the acceleration 𝑎 of the moving object is determined by the
formula 𝐹 = 𝑚𝑎.
ACCEPTANCE REGION: The set of values for a random variable in hypothesis
testing such that the null hypothesis is not rejected.

ACCEPTANCE REJECTION ALGORITHM: An algorithm for generating random numbers

from some probability distribution, 𝑓(𝑥), by first generating a random number from
some other distribution, 𝑔(𝑥), where 𝑓 and 𝑔 are related by 𝑓(𝑥) ≤ 𝑔(𝑥) for all 𝑥 with 𝑘
a constant. The algorithm works as follows:

1. let 𝑟 be a random number from 𝑔(𝑥);

2. let 𝑠 be a random number from a uniform distribution on the interval (0,1);
3. calculate 𝑐 = 𝑘𝑠𝑔(𝑟);
4. if 𝑐 > 𝑓 (𝑟) reject 𝑟 and return to the first step; if 𝑐 ≤ 𝑓 (𝑟) accept 𝑟 as a random
number from 𝑓.

ACCEPTANCE SAMPLING: A type of quality control procedure in which a sample is

taken from a collection or batch of items, and the decision to accept the batch as
satisfactory, or reject them as unsatisfactory, is based on the proportion of defective
items in the sample. It is a method of quality control where a sample is taken and a
decision whether to accept the population is made on the basis of the quality of the
sample. The simplest method is to have a straight accept/reject criterion, but a more
classy approach is to take another sample if the evidence from the existing sample, or a
set of samples, is not evidently indicating whether the group should be accepted or
rejected. One of the main benefits of this approach is reducing the cost of taking samples
to satisfy quality control criteria.

ACCUMULATION POINT: An accumulation point is a point which is the limit of a

sequence, also called a limit point. For some maps, periodic orbits give way to chaotic
ones beyond a point known as the accumulation point.
ACCURACY: Accuracy is defined as a measure of the exactness of a numerical quantity,
frequently given to n significant figures (where the proportional accuracy is the
significant aspect) or n decimal places (where the absolute accuracy is more
significant).
ACCURATE (CORRECT) TO n DECIMAL PLACES: Rounding a number with the accuracy
specified by the number of decimal places given in the rounded value. So e = 2.71828 …
= 2.718 to three decimal places and = 2.72 to two decimal places.
ACTIVE CONSTRAINT: An inequality such as 3𝑦 + 5𝑥 ≥ 21 is said to be active at a
point on the boundary, i.e. where equality holds, for example (2, 3) 𝑎𝑛𝑑 (7, 0).
ACUTE ANGLE: An angle whose measure is less than 90 degrees.
ACYCLIC AND CYCLIC MOTION: When the region occupied by the fluid moving
irrotationally is simply connected, the velocity potential will be single-valued. Thus the
velocity potential becomes

∅B = − rdr
OAB

For all paths from O to B. Thus the motion in which the velocity potential is single
valued is called acyclic. In simply connected region the only possible irrotational motion
is acylic. When the velocity potential is not single- valued, the motion is said to be cyclic
i.e., it is not possible to assign to every point of the original region a unique and definite
value of ∅.

ACYCLIC GRAPH: A graph is said to be acyclic if it contains no circuits.

ADAPTIVE CLUSTER SAMPLING: A procedure in which an initial set of subjects is
selected by some sampling procedure and, whenever the variable of interest of a
selected subject satisfies a given criterion, additional subjects in the neighbourhood of
that subject are added to the sample.

ADDEND: A quantity to be added to another, also called a summand. For example, in

the expression 𝑎 + 𝑏 + 𝑐, 𝑎, 𝑏, and 𝑐 are all addends. The first of several addends, or “the
one to which the others are added” a in this example , is sometimes called the augend.
One of the numbers that are being added together in a sum. For example, when we write
4 + 5 = 9, 4 𝑎𝑛𝑑 5 are called the addends.
ADDITION: The operation, or process, of calculating the sum of two numbers or
quantities.
ADDITION (OF COMPLEX NUMBERS): Let the complex numbers 𝑧1 and 𝑧2 , where
𝑧1 = 𝑎 + 𝑖𝑏 and 𝑧2 = 𝑐 + 𝑖𝑑 be represented by the points 𝑃 and 𝑄 in the complex plane.
Then 𝑧1 + 𝑧2 = (𝑎 + 𝑐) + 𝑖(𝑏 + 𝑑), and 𝑧1 + 𝑧2 is represented in the complex plane
by the point 𝑅 such that 𝑂𝑃𝑅𝑄 is a parallelogram.

ADDITION (OF LINEAR MAPS): We define a map 𝑇1 + 𝑇2 : 𝑈 → 𝑉 by the rule

( 𝑇1 + 𝑇2 )(𝑢) = 𝑇1 𝑢 + 𝑇2 (𝑢) 𝑓𝑜𝑟 𝑢 𝜖 𝑈

ADDITION (OF MATRICES): Let 𝑨 and 𝑩 be 𝑚 × 𝑛 matrices, with 𝑨 = [𝑎𝑖𝑗 ] and

𝑩 = [𝑏𝑖𝑗 ]. The operation of addition is defined by taking the sum 𝑨 + 𝑩 to be the
𝑚 × 𝑛 matrix 𝑪, where 𝑪 = [𝑐𝑖𝑗 ] and 𝑐𝑖𝑗 = 𝑎𝑖𝑗 + 𝑏𝑖𝑗 . The sum 𝑨 + 𝑩 is not defined if
𝑨 and 𝑩 are not of the same order. This operation + of addition on the set of all m × n
matrices is associative and commutative. Let 𝐴 and 𝐵 be two matrices of the same type
𝑚 × 𝑛. The their sum (to be denoted by 𝐴 + 𝐵) is defined to be the matrix of the type
𝑚 × 𝑛 obtained by the adding the corresponding elements of 𝐴 and 𝐵. Thus if

𝐴 = 𝑎𝑖𝑗 and 𝐵 = 𝑏𝑖𝑗 ,then 𝐴 + 𝐵 = 𝑎𝑖𝑗 + 𝑏𝑖𝑗 .

𝑚 ×𝑛 𝑚 ×𝑛 𝑚 ×𝑛

Note that A + B is also a matrix of the type 𝑚 × 𝑛. More clearly we can say that:

a11 a12 … a1n b11 b12 … b1n

a21 a22 … a2n b21 b22 … b2n
A= … … … … and B =
… … … …
am1 am2 … amn m×n b m1 bm2 … bmn m×n

Then
 a11 + b11 a12 + b12 … a1n + b1n
a + b21 a22 + b22 … a2n + b2n
A + B = 21
……… ……… … ……
am1 + bm1 am2 + bm2 … amn + bmn m×n.

For example, if

3 2 −1 1 −2 7
A= and B = ; then
4 −3 1 2×3 3 2 −1 2×3

3+1 2 − 2 −1 + 7 4 0 6
A+B= =
4+3 −3 + 2 1 − 1 7 −1 0 2×3.

It should be noted that addition is defined only for matrices which are of the same size.
If two matrices A and B are of the same size, they are said to be conformable for
addition. If the matrices A and BA are not of the same size, we cannot find their sum.

ADDITION (OF VECTORS): Given vectors 𝒂 and 𝒃, let 𝑂𝐴 and 𝑂𝐵 be directed line-
segments that represent 𝒂 and 𝒃, with the same initial point O. The sum of 𝑂𝐴 and 𝑂𝐵 is
the directed line-segment 𝑂𝐶 , where OACB is a parallelogram, and the sum 𝒂 + 𝒃 is
defined to be the vector c represented by 𝑂𝐶 .
This is called the parallelogram law. Alternatively, the sum of vectors 𝒂 and 𝒃 can be
defined by representing a by a directed line-segment 𝑂𝑃 and 𝒃 by 𝑃𝑄 where the final
point of the first directed line-segment is the initial point of the second. Then 𝒂 + 𝒃 is
the vector represented by 𝑂𝑄. This is called the triangle law. Addition of vectors has the
following properties, which hold for all a, b and c:
(i) 𝒂 + 𝒃 = 𝒃 + 𝒂, the commutative law.
(ii) 𝒂 + (𝒃 + 𝒄) = (𝒂 + 𝒃) + 𝒄, the associative law.
(iii) 𝒂 + 𝟎 = 𝟎 + 𝒂 = 𝒂, where 0 is the zero vector.
(iv) 𝒂 + (−𝒂) = (−𝒂) + 𝒂 = 𝟎, where −𝒂 is the negative of 𝒂.
ADDITION LAW: In the theory of probability, If A and B are two events then the addition
law states that the
𝑃(𝐴 ∪ 𝐵) = 𝑃(𝐴) + 𝑃(𝐵) − 𝑃(𝐴 ∩ 𝐵).
In the special case, where A and B are mutually exclusive events this reduces to
𝑃(𝐴 ∪ 𝐵) = 𝑃(𝐴) + 𝑃(𝐵).
ADDITIVE FUNCTION: A function 𝑓(𝑥) over a domain 𝐷 is said to be an additive
functiuon if
𝑓 𝑥 + 𝑦 = 𝑓 𝑥 + 𝑓 𝑦 ∀𝑥, 𝑦 ∇ 𝐷.
For example, 𝑓(𝑥) = 𝑛𝑥 is an additive function since
𝑓 𝑥 + 𝑦 = 𝑛 𝑥 + 𝑦 = 𝑛𝑥 + 𝑛𝑦 = 𝑓(𝑥) + 𝑓(𝑦)
ADDITIVE GROUP: A group G with the operation + (addition) may be called an additive
group.
ADDITIVE IDENTITY: The number zero is known as the additive identity element,
because it satisfies the property that the addition of zero does not change a number:
𝑎 + 0 = 𝑎 = 0 + 𝑎 ∀ 𝑎 ∇ 𝑅..
ADDITIVE INVERSE: The sum of a number and its additive inverse is zero. The additive
inverse of 𝑎 (written as −𝑎) is also called the negative or the opposite of 𝑎: 𝑎 + −𝑎 =
0 = −𝑎 + 𝑎 ∀ 𝑎 ∇ 𝑅. For example, -6 is the additive inverse of 6 and -7 is the additive
inverse of 7.
ADHERENT POINT: Adherent point of a set A is a point of the set A in a topological space
that belongs to the closure of a set A. In other words, a point 𝑝 ∇ 𝑆 is said to be an
adherent point of 𝑆 if every nbd of 𝑝 contains a point of 𝑆. The set of all adherent points
of the set 𝑆 is called adherence of 𝑆 and is denoted by 𝐴𝑑𝑕 𝐴.
ADHERENT POINT OF A FILTER: Let (𝐸, 𝑇) be a topological space and 𝐹 a filter on 𝐸. A
point 𝑥 ∇ 𝐸 is said to be adherent point of 𝐹 if it is adherent to every 𝐴 ∇ 𝐹, i.e. if every
𝑋 ∇ 𝐵(𝑥) meets every 𝐴 ∇ 𝐹. Consider in 𝑅 the set of points 1/𝑛, 1 − 1/𝑛 (where
𝑛 ∇ 𝑁) and the point 2. Let 𝐹 be the filter consisting of the complements of finite
subsets of this set. 𝐹 has the points 0 and 1 as adherent points.
AD INFINITUM: Literally meaning "to infinity" in Latin. It is a phrase used in
mathematics (and other fields) as a more elegant way of what is more commonly
expressed as "... and so on and so forth".

ADJ: It is the abbreviation for adjoint of a matrix.

ADJACENCY MATRIX: For a simple graph G, with n vertices 𝑣1 , 𝑣2 , … , 𝑣𝑛 , the adjacency
matrix A is the 𝑛 × 𝑛 matrix [𝑎𝑖𝑗 ] with 𝑎𝑖𝑗 = 1, if 𝑣𝑖 is joined to 𝑣𝑗 , and 𝑎𝑖𝑗 = 0,
otherwise. The matrix A is symmetric and the diagonal entries are zero. The number of
ones in any row (or column) is equal to the degree of the corresponding vertex. An
example of a graph and its adjacency matrix A is shown below.

ADJACENCY RELATION: The set E of edges of a graph (V, E), being a set of unordered
pairs of elements of V, constitutes a relation on V. Formally, an adjacency relation is any
relation which is irreflexive and symmetric.

ADJACENT ANGLES: Two angles that share a ray, thereby being directly next to each
other
ADJACENT EDGES: Adjacent edges are a pair of edges in a graph joined by a common
vertex.
ADJACENT FRACTION: Two fractions are said to be adjacent if their difference has a
unit numerator. For example, I/3 and 1/4 are adjacent since l/3 - l/4 = l/12, but l/2 and
l/5 are not since l/2 - l/5 = 3/10. Adjacent fractions can be adjacent in a Farey sequence.

ADJACENT VERTICES: Adjacent vertices are a pair of vertices in a graph joined by a

common edge.
ADJOINT CURVE: A curve which has at least multiplicity 𝑝𝑖 − 1 at each point where a
given curve (having only ordinary singular points and cusps) has a multiplicity Q is
called the adjoint to the given curve. When the adjoint curve is of order 𝑛 − 3, it is
called a special adjoint curve.

ADJOINT OF A SQUARE MATRIX: Let A = aij be any n × n matrix. The transpose 𝐵′

n×n

of the matrix B = Aij , where Aij denotes the cofactor of the element aij in the
n×n

determinant A , is called the adjoint of the matrix A and is denoted by ht symbol Adj. A.
If A be any n-rowed square matrix, then Adj A A = A Adj A = A In , where 𝐼𝑛 is the n-
rowed unit matrix.

ADJOINT OPERATOR: Let 𝐻 and 𝐾 be Hilbert Spaces and 𝑇 ∇ 𝐵(𝐻, 𝐾). Then there exists
operator 𝑇 ∗∇ 𝐵(𝐾, 𝐻) such that ⟨ 𝑇𝑕, 𝑘 ⟩𝐾 = ⟨ 𝑕, 𝑇 ∗ 𝑘 ⟩𝐻 for all 𝑕 ∇ 𝐻, 𝑘 ∇ 𝐾. Such
𝑇 ∗ is called the adjoint operator of 𝑇. Also 𝑇 ∗∗ = 𝑇 𝑎𝑛𝑑 ||𝑇 ∗ || = ||𝑇||.

 For operators 𝑇1 and 𝑇2 , (𝑇1 𝑇2 )*= 𝑇2 *𝑇1 *, (𝑇1 +𝑇2 )*= 𝑇1 *+𝑇2 * and λ T)*=λT*.
 If A is an operator on a Hilbert space H then (𝑘𝑒𝑟𝐴)⊥ = 𝐼𝑚 𝐴∗ .

ADJUNCTION: If 𝑎 is an element of a field 𝑓 over the prime field 𝑃, then the set of all
rational functions of a with coefficients in 𝑃 is a field derived from 𝑃 by adjunction of 𝑎.

ADMISSIBLE: A string or word is said to be admissible if that word appears in a given

sequence. For example, in the sequence aabaabaabaabaab . . ., a, aa, baab are all
admissible, but bb is inadmissible.

ADMISSIBLE HYPOTHESIS: Any hypothesis that has not been ruled out, that is, a
hypothesis that may possibly be true.
a.e.: It is the abbreviation for almost everywhere
AEROFOIL: The aerofoil has a profile if fish type. It is used in modern airplanes. Such an
aerofoil has a blunt leading edge and a sharp trailing edge. The projection of the profile
on the double tangent is the chord. The ratio of the span to the chord is the aspect ratio.
The locus of the point midway between the points in which an ordinate perpendicular
to the chord meets the profile is known as the camber line of a profile. The camber is the
ratio of the maximum ordinate of the camber line to the chord.

The theory of the flow round such an aerofoil is made on the following assumptions:

1. The air behaves as an incompressible inviscid fluid,

 2. The aerofoil is as cylinder whose cross section is a curve of the above form and
the flow is two-dimensional irrotational cyclic motion.

AFFINE GEOMETRY: Affine Geometry is not concerned with the notions

of circle, angle and distance. It's a known dictum that in Affine Geometry all triangles are
the same. In this context, the word affine was first used by Euler (affinis). In modern
parlance, Affine Geometry is a study of properties of geometric objects that remain
invariant under affine transformations (mappings). Affine transformations preserve
collinearity of points: if three points belong to the same straight line, their images under
affine transformations also belong to the same line and, in addition, the middle point
remains between the other two points. As further examples, under affine
transformations

 parallel lines remain parallel,

 concurrent lines remain concurrent (images of intersecting lines intersect),
 the ratio of length of line segments of a given line remains constant,
 the ratio of areas of two triangles remains constant (and hence the ratio of any
areas remain constant),
 ellipses remain ellipses and the same is true for parabolas and hyperbolas.
 barycenters of triangles (and other shapes) map into the corresponding
barycenters.

AFFINE SPACES: An affine space 𝐴 is constructed as follows: Let 𝑉 be a vector space

over a field 𝐾, and let 𝐴 be a nonempty set. For any vector 𝑎 ∇ 𝑉 and any element 𝑝 of 𝐴,
suppose that an addition 𝑝 + 𝑎 ∇ 𝐴 is defined so as to satisfy the following three
conditions:
(i) 𝑝 + 0 = 𝑝 (0 being a zero vector);
(ii) (𝑝 + 𝑎) + 𝑏 = 𝑝 + (𝑎 + 𝑏) (𝑎, 𝑏 ∇ 𝑉); and
(iii) for any 𝑞 ∇ 𝐴 there exists a unique vector 𝑎 ∇ 𝑉 such that 𝑞 = 𝑝 + 𝑎.
Condition (i) follows from (ii) and (iii). Then we call 𝐴 an affine space, 𝑉 the standard
vector space of 𝐴, and 𝐾 the coefficient field of 𝐴. Each element of 𝐴 is called a point.
AFFINE TRANSFORMATION: A transformation which involves any combination
of translations, reflections, stretches, shrinks, or rotations. Collinearity and concurrency
are invariant under affine transformations.
AGGREGATION: The process of representing the sum of mathematical terms as a single
(mathematical) term.

AGREE: If 𝑓(𝑥) and 𝑔(𝑥) are two functions defined on a set S, and 𝑓(𝑥) = 𝑔(𝑥) for all
𝑥 ∇ 𝑆, then we say that 𝑓 and 𝑔 agree on the set 𝑆.
Ai: It is the abbreviation for Airy function.
AHLFORS-BERS THEOREM: The Riemann’s moduli space gives the solution to
Riemann’s moduli problem, which requires an analytic parameterization of the compact
Riemann surfaces in a fixed homeomorphism.

AIRLINE PROBLEM: Given n cities, it is desired to design a network of flights between

the cities satisfying the following conditions. For each city there are direct flights to and
from exactly the same number of other cities. Between any two cities there is exactly
one flight with at most one stopover. For what values of n do such networks exist?

AITKEN’S METHOD (NUMERICAL METHODS): If an iterative formula 𝑥𝑟+1 = 𝑓(𝑥𝑟 ) is to

be used to solve an equation, Aitken’s method of accelerating convergence uses the
initial value. This can then be used as a new starting point from which to repeat the
process until the required accuracy has been reached. Let us suppose that 𝑥0 , 𝑥1 , 𝑥2 are
the initial value and the first two iterations and
𝛥𝑥𝑟 = 𝑥𝑟+1 − 𝑥𝑟 , 𝛥2 𝑥𝑟 = 𝛥𝑥𝑟+1 − 𝛥𝑥𝑟
are the forward differences then
𝛥𝑥 2 2
𝑥4 = 𝑥3 − .
𝛥2 𝑥 1

More generally this will be expressed as

𝛥𝑥𝑟−1 2
𝑥𝑟+1 = 𝑥𝑟 − .
𝛥2 𝑥𝑟−2
ALAOGLU THEOREM: The unit ball in 𝑋 ∗ is weak compact.
ALEPH: Aleph is a symbol that is used to denote any infinite cardinal number and
usually denoted by the Hebrew letter ‫א‬.
ALEPH-NULL: Aleph-null is the smallest infinite cardinal number. It is the cardinality of
any set which can be put in one-to-one correspondence with the set of natural numbers.
Such sets are also known as countable or denumerable sets. One of the apparent
paradoxes in number theory is that the set of rational numbers between 0 and 1, the set
of all rational numbers, and the set of natural numbers all have the same cardinality.
The symbol ‫א‬0 is used.
ALEXANDROFF COMPACTIFICATION: The collection consisting of all open subsets of 𝑋
and all complements in 𝑋 ∞ of closed compact subsets of 𝑋 is the finest topology on 𝑋 ∞
such that 𝑋 ∞ is compact and contains 𝑋 as a subspace. If 𝑋 is not compact then 𝑋 is
dense in 𝑋 ∞ , and 𝑋 ∞ is called the Alexandroff compactification of 𝑋.

ALEXANDROV SPACE: A generalization of Riemannian manifolds with upper, lower or

integral curvature bounds (the last one works only in dimension 2)
ALEXANDROV TOPOLOGY: The topology of a space 𝑋 is an Alexandrov topology (or
is finitely generated) if arbitrary intersections of open sets in 𝑋 are open, or
equivalently, if arbitrary unions of closed sets are closed, or, again equivalently, if the
open sets are the upper sets of a poset.
ALGEBRA: Algebra is the branch of Mathematics that studies the properties of
operations carried out on sets of numbers. It started with Alkwarizmi’s 820 Al Gebr
W’ali Muquabala, the origin of the word “algebra.” It was the first mathematical book
written in Arabic. Its content was essentially a variety of methods of solving algebraic
equations. Al gebr means “transposition of negative terms on one side of the equation to
the other side and changing their signs.” and W’ali muquabala means “simplification of
the equation by gathering similar terms.” Algebra is a generalization of arithmetic in
which symbols, usually letters, are used to stand for numbers. It is the area of
mathematics related to the general properties of arithmetic. Relationships can be
summarized by using variables, usually denoted by letters x, y, n, … to stand for
unknown quantities, whose value(s) may be determined by solving the resulting
equations. The structure of algebra is based upon axioms (or postulates), which are
statements that are assumed to be true. Some algebraic axioms include the transitive
axiom and the associative axiom of addition i.e.:
 𝑎 = 𝑏 𝑎𝑛𝑑 𝑏 = 𝑐 ⇒ 𝑎 = 𝑐.
 𝑎+ 𝑏 + 𝑐= 𝑎+ 𝑏+ 𝑐 .
These axioms are used to prove theorems about the properties of operations on real
numbers.
ALGEBRAICALLY CLOSED: A field 𝐹 is said to be algebraically closed if every polynomial
with coefficients in 𝐹 has a solution in 𝐹. Equivalently, every polynomial with
coefficients in 𝐹 can be written as a product of linear polynomials.
ALGEBRAIC CLOSURE: It is the procedure of the extension of a given set to include all
the roots of polynomials with coefficients in the given set. The smallest algebraically
closed set of numbers is C, the set of complex numbers, since the very simple equation
𝑥 2 + 1 = 0 has a complex solution.
ALGEBRAIC FIELD EXTENSION: Let 𝐿: 𝐾 be a field extension, and let 𝛼 be an element of
𝐿. If there exists some non-zero polynomial 𝑓 ∇ 𝐾[𝑥] with coefficients in 𝐾 such that
𝑓(𝛼) = 0, then α is said to be algebraic over 𝐾; otherwise 𝛼 is said to be transcendental
over 𝐾. A field extension 𝐿: 𝐾 is said to be algebraic if every element of 𝐿 is algebraic
over 𝐾.
ALGEBRAIC FUNCTION: An algebraic function is a function which can be expressed as a
root of an equation of the form
𝜆𝑛 + 𝑏𝑛–1 𝜆𝑛–1 + ··· + 𝑏1 𝜆 + 𝑏0 = 0
Where 𝑏𝑛–1 , − − − 𝑏1 , 𝑏0 are rational functions of 𝑥.
ALGEBRAIC GEOMETRY: Algebraic geometry is the branch of mathematics that deals
with algebraic varieties, that is, point sets defined by several algebraic equations in a
space of any dimension or those derived from these sets by means of certain
constructions, It may also be considered to be a theory of the field of algebraic functions
in several variables in geometric language, and it is closely related to the theories of
complex analytic manifolds, commutative algebra, and homological algebra. It also has
an important connection with number theory through the theories of automorphic
functions, Diophantine equations, and zeta functions.
ALGEBRAIC INTEGER: A complex number 𝑧 is said to be an algebraic integer if it is the
root of some monic polynomial with integer coefficients. It follows from the above
definition that a complex number is an algebraic integer if and only if it is integral over
the ring of integers. It follows from the relevant definitions that a complex number is an
algebraic number if and only if it is integral over the field 𝑄 of rational numbers. Also a
complex number is an algebraic integer if and only if it is integral over the ring 𝑍 of
(rational) integers. In algebraic number theory, elements of the ring 𝑍 are often referred
to as rational integers to distinguish them from algebraic integers. The set of all
algebraic integers constitutes a subring of the field 𝐶 of complex numbers that is
integrally closed in 𝐶. An algebraic number is an algebraic integer if and only if the
coefficients of its minimum polynomial over the field 𝑄 of rational numbers are rational
integers.
ALGEBRAIC NOTATION: This avoids the initial use of components, and is distinguished
by the explicit use of the tensor product symbol.
ALGEBRAIC NUMBER: Algebraic number is a real number that is the root of some
polynomial equation with integer coefficients. All rational numbers are algebraic, since
𝑎/𝑏 is the root of the equation 𝑏𝑥 − 𝑎 = 0. Some irrational numbers are algebraic; for
example, 3 is the root of the equation 𝑥 2 − 3 = 0. An irrational number that is not
algebraic such as π, e etc. is called a transcendental number.
ALGEBRAIC NUMBER FIELD: A complex number that satisfies an algebraic equation with
rational integral coefficients is said to be an algebraic number. If the coefficient of the
term of highest degree of the equation is 1, this algebraic number is said to be an
algebraic integer. The set 𝐴 of ah algebraic numbers is a field which is the algebraic
closure of the rational number field 𝑄 in the complex number field 𝐶. The set 𝐼 of all
algebraic integers is an integral domain which contains the integral domain 𝒁 of all the
rational integers. The field of quotients of 𝐼 is 𝐴. An extension field 𝑘 of 𝑄 of finite
degree, which we shall always suppose to be contained in C is said to be an algebraic
number field of finite degree, and 𝑘 is a subfield of 𝐴. Thus An algebraic number field is a
subfield of the field 𝐶 of complex numbers that is a finite-dimensional vector space over
the field 𝑄 of rational numbers.
ALGEBRAIC SET: A subset of 𝑛-dimensional affine space 𝐴𝑛 is said to be an algebraic set
if it is of the form {(𝑥1 , 𝑥2 , . . . , 𝑥𝑛 ) ∇ 𝐴𝑛 ∶ 𝑓(𝑥1 , 𝑥2 , . . . , 𝑥𝑛 ) = 0 for all 𝑓 ∇ 𝑆} for some
subset 𝑆 of the polynomial ring 𝐾[𝑋1 , 𝑋2 , . . . , 𝑋𝑛 ].
ALGEBRAIC STRUCTURE: It is the term used to describe an abstract concept defined as
consisting of certain elements with operations satisfying given axioms. Thus, a group
(ring or field) is an algebraic structure. The purpose of the definition is to recognize
similarities that appear in different contexts within mathematics and to encapsulate
these by means of a set of axioms.
ALGEBRAIC SYSTEM: It is a set together with the operations and relations defined on
that set.
ALGEBRA (LINEAR ALGEBRA): A linear space 𝑆 over a field 𝐹 is called an algebra, (or a
linear algebra) if a composition (.) is defined in 𝑆 such that
 (i) 𝑓. 𝑔 . 𝑕 = 𝑓. 𝑔. 𝑕 𝑓𝑜𝑟 𝑎𝑙𝑙 𝑓, 𝑔, 𝑕 ∇ 𝑆,
(ii) 𝑓. 𝑔 + 𝑕 = 𝑓. 𝑔 + 𝑓. 𝑕 and 𝑓 + 𝑔 . 𝑕 + 𝑔. 𝑕, for all 𝑓, 𝑔, 𝑕, ∇ 𝑆, and
(iii) 𝑎 𝑓. 𝑔 = 𝑎𝑓 . 𝑔 = 𝑓. 𝑎𝑔 , for all 𝑓, 𝑔, 𝑕 ∇ 𝑆, and 𝑎 ∇ 𝐹.

ALGEBRA OF MEASURABILITY: Let 𝑓, 𝑔: 𝑋 → ℝ be measurable. Then the functions

𝑎𝑓, 𝑓 + 𝑔, 𝑓𝑔, 𝑚𝑎𝑥(𝑓, 𝑔) 𝑎𝑛𝑑 𝑚𝑖𝑛(𝑓, 𝑔) are all measurable. That is measurable functions
form algebra and this algebra is closed under convergence a.e.

We say that sequence (𝑓𝑛 ) of functions converges

𝑠𝑢𝑝
1. uniformly to f (notated 𝑓𝑛 ⇒ 𝑓) if 𝑓 𝑥 − 𝑓(𝑥) → 0
𝑥∇𝑋 𝑛
𝑎. 𝑒
2. almost everywhere to 𝑓 (notated 𝑓𝑛 𝑓) if 𝑓𝑛 ⇒ 𝑓 for all 𝑥 ∇ 𝑋\𝐴, 𝜇 𝐴 = 0.
→
𝜇
3. in measure µ to f (notated 𝑓𝑛 𝑓) if for all 𝜀 > 0, 𝜇 {𝑥 ∇ 𝑋: 𝑓𝑛 𝑥 − 𝑓(𝑥) >
→
𝜀} → 0.

Clearly uniform convergence implies both convergences a.e and in measure.

On finite measures convergence a.e. implies convergence in measure.

ALGEBRA OF SETS: The set of all subsets of a universal set E is closed under the binary
operations ∪ union and ∩ intersection and the unary operation complementation .
The following are some of the properties, or laws, that hold for subsets A, B and C of E:
(i) 𝐴 ∪ (𝐵 ∪ 𝐶) = (𝐴 ∪ 𝐵) ∪ 𝐶 and 𝐴 ∩ (𝐵 ∩ 𝐶) = (𝐴 ∩ 𝐵) ∩ 𝐶 (Associative
property)
(ii) 𝐴 ∪ 𝐵 = 𝐵 ∪ 𝐴 and 𝐴 ∩ 𝐵 = 𝐵 ∩ 𝐴 (commutative property)
(iii) 𝐴 ∪ Ø = 𝐴 and 𝐴 ∩ Ø = Ø, where Ø is the empty set.
(iv) 𝐴 ∪ 𝐸 = 𝐸 and 𝐴 ∩ 𝐸 = 𝐴.
(v) 𝐴 ∪ 𝐴 = 𝐴 and 𝐴 ∩ 𝐴 = 𝐴.
(vi) 𝐴 ∩ (𝐵 ∪ 𝐶) = (𝐴 ∩ 𝐵) ∪ (𝐴 ∩ 𝐶)
and 𝐴 ∪ (𝐵 ∩ 𝐶) = (𝐴 ∪ 𝐵) ∩ (𝐴 ∪ 𝐶) (distributive property)
(vii) 𝐴 ∪ 𝐴′ = 𝐸 and 𝐴 ∩ 𝐴′ = Ø.
(viii) 𝐸′ = Ø and Ø′ = 𝐸.
(ix) (𝐴′)′ = 𝐴.
(x) (𝐴 ∪ 𝐵)′ = 𝐴′ ∩ 𝐵′ and (𝐴 ∩ 𝐵)′ = 𝐴′ ∪ 𝐵′ De Morgan’s laws
The application of these laws to subsets of E is known as the algebra of sets.
ALGEBRA WITH IDENTITY: An algebra A is called an algebra with identity if there exists
a non-zero element in A, denoted by 𝑒 and called the identity element or the identity
such that
𝑒𝑓 = 𝑓𝑒 = 𝑓 𝑓𝑜𝑟 𝑎𝑙𝑙 𝑓𝜖 𝐴
Obviously the identify in an algebra is unique.
ALGORITHM: An algorithm is a chain of instructions that tell how to carry out a task. An
algorithm must be specified exactly, so that there can be no doubt about what to do
next, and it must have a finite number of steps.
ALIQUOT PART : Aliquot part is another name for a proper divisor, i.e. any divisor of a
given number other than the number itself. A prime number has only one aliquot part -
the number 1. 1, 2, 3, 4, 6 are all aliquot parts of 12. The number 1 does not have aliquot
parts.
AL-KHWARIZMI: Muhammad Ibn Musa Al-Khwarizmi (780 AD to 850 AD) was a
Muslim mathematician who introduced modern numerals (the Hindu-arabic numerals)
to Europe in his book titled Kitab -al-jabr-walli-muqabalah that provided the source for
the term algebra. His name is the source for the term algorithm.
ALMOST CONTINUOUS FUNCTION: A function 𝑓 ∶ 𝑋 → 𝑌 is almost continuous,
𝑓 ∇ 𝐴𝐶(𝑋, 𝑌 ), provided each open subset of 𝑋 × 𝑌 containing the graph of 𝑓 also
contains the graph of a continuous function from 𝑋 to 𝑌 . This property was defined as a
generalization of functions having the fixed-point property. It is easy to see that if every
function in 𝐶(𝑋, 𝑋) has the fixed-point property, then so does every 𝑓 ∇ 𝐴𝐶(𝑋, 𝑋).

ALMOST EVERYWHERE: If we have a property 𝑃(𝑥) which is true except possibly for
𝑥 ∇ 𝐴 and µ(𝐴) = 0, we say 𝑃(𝑥) is almost everywhere or a.e. t is Therefore, it is a
property holding for all values except on a set of zero measure. The most remarkable
example of a set of zero measure is the set of rational numbers, so if

1, when 𝑥 is rational
𝑓 𝑥 = ,
0, when 𝑥 is irrational
then we can say that 𝑓(𝑥) = 0 almost everywhere.
ALMOST PERIODIC FUNCTIONS: The theory of almost periodic functions was originated
by H. Bohr in 1924 as a result of his study of Dirichlet series. The theory provides a
method of studying a wide class of trigonometric series of general type.
Let 𝑓(𝑥) be a complex-valued continuous function defined for all real values of 𝑥. A
number 𝑻 is called a translation number of 𝑓(𝑥) belonging to 𝜖 > 0 if
 𝑠𝑢𝑝
𝑓(𝑥 + 𝑇) − 𝑓(𝑥) <∇
−∞ < 𝑥 < ∞
If for any ∇ > 0 there exists a number 𝑙(∇) > 0 such that any interval of length 𝑙(∇)
contains a translation number belonging to ∇, then 𝒇(𝒙) is called almost periodic in the
sense of Bohr. We denote by 𝑩 the set of all almost periodic functions in the sense of
Bohr.
ALPHABETS AND WORDS: Let 𝐴 be a finite set. We refer to this set 𝐴 as an alphabet and
we refer to the elements of 𝐴 as letters. (For example, the set 𝐴 might be the set of
letters in the English language, or in any other world language, or the set {0, 1} of binary
digits, or the set {0, 1, 2, 3, 4, 5, 6, 7, 8, 9} of decimal digits.) For any natural number 𝑛, we
define a word of length 𝑛 over the alphabet 𝐴 to be a string of the form 𝑎1 𝑎2 . . . 𝑎𝑛 in
which 𝑎𝑖 ∇ 𝐴 for 𝑖 = 1, 2, . . . , 𝑛. We shall denote by 𝐴𝑛 the set of all words of length 𝑛
over the alphabet 𝐴. In particular, 𝐴1 = 𝐴.
Alt: alternating group (Alt(n) is also written as An)
ALTERNATING GROUP, An: It is the subgroup of the symmetric group 𝑆𝑛 and it contains
all the even permutations of 𝑛 objects. It has order 𝑛!/2. For 𝑛 > 4, it is the only proper
normal subgroup of 𝑆𝑛 apart from the empty set. The alternating group is a simple
group.
ALTERNATING MAP: A multilinear map 𝜙: 𝑉 𝑘 = 𝑉 × · · · × 𝑉 → 𝑈 where 𝑘 > 1, is
said to be alternating if for all 𝑣1 , . . . , 𝑣𝑘 ∇ 𝑉 the value 𝜙(𝑣1 , . . . , 𝑣𝑘 ) changes sign
whenever two of the vectors 𝑣1 , . . . , 𝑣𝑘 are interchanged, that is
𝜙(𝑣1 , . . . , 𝑣𝑖 , . . . , 𝑣𝑗 , . . . , 𝑣𝑘 ) = −𝜙(𝑣1 , . . . , 𝑣𝑗 , . . . , 𝑣𝑖 , . . . , 𝑣𝑘 ).
ALTERNATING SERIES: An alternating series is a series in which every term has the
opposite sign from the preceding term. For example, 1 − 𝑥 + 𝑥 2 − 𝑥 3 + − − − is an
alternating series.
ALTERNATIVE HYPOTHESIS: The model considered to be the case if the null
hypothesis in considered to be rejected (not hold). There is some debate as to whether
an alternative hypothesis increases the subjectiveness of hypothesis testing, since the
alternative hypothesis does not necessarily have to be the negation of the null
hypothesis. The issue is whether it is appropriate to make presumptions of what
the model in consideration cannot be (models not covered by the null or alternative
hypothesis but not considered) in hypothesis testing. Setting the alternative hypothesis
to be the negation of the null hypothesis amounts to not using alternative hypothesis as
some statistician would advocate. The alternative hypothesis is the hypothesis that
states, “The null hypothesis is false.”

AMBIGUOUS: Having more than one possible way, meaning , value, or solution.
AMICABLE NUMBERS: A pair of numbers with the property that each is equal to the sum
of the positive proper divisors of the other. Here proper divisor means that the number
is not included as one of its own divisors. For example, 220 and 284 are amicable
numbers because the positive divisors of 220 are 1,
2, 4, 5, 10, 11, 20, 22, 44, 55 and 110, whose sum is 284, and the positive divisors of 284
are 1, 2, 4, 71 and 142, whose sum is 220. More than 600 pairs are now known.
AMOUNT: The total value of an investment, including the principal and the interest
AMPLITUDE: Half the difference between the maximum and the minimum values of a
periodic function. The amplitude of a periodic function is half of the difference between
the largest possible value of the function and the smallest possible value. For example,
for the function 𝑦 = 𝑠𝑖𝑛 𝑥, the largest possible value of 𝑦 is 1 and the smallest possible
1−(−1)
value of 𝑦 is −1, so the amplitude is = 1.
2

An : The alternating group of degree 𝑛 and order 𝑛!/2. That is, a group having 𝑛!/
2 members which act on 𝑛 elements of a set. It is a group that consists of members
which are even permutations (compositions of even number of permutations) only.

ANALOG COMPUTATION: The term analog computation is a generic term describing

various techniques of computation employing diagrams, or physical systems whose
equations are similar to the mathematical problems in question. The history of analog
computation is probably as old as that of digital computation; for example, the ancient
Greeks tried to solve cubic equations using diagrams, and the astrolabe widely used by
astronomers through the medieval period is also a kind of analog computer. Soon after
the discovery of logarithms, the slide rule was invented. In the 18th Century, the
planimeter, used to measure plane areas, was introduced, and in the 19th Century, the
nomogram appeared. In the first half of 20th Century, a large electronic analog computer
was developed, thus predating the first practical digital computer. However, analog
computation has an essential defect, namely, the limitation of accuracy. Today it is
useful for simple calculation, but is rapidly becoming less important as the development
of digital computers, including pocket calculators, advances.
ANALYSIS: Analysis (Real Analysis, Mathematical Analysis or Complex Analysis) is the
branch of mathematics that studies limits and convergence; calculus is a part of analysis.
ANALYSIS OF VARIANCE: The sum of squared deviations from the mean of a whole
sample is split into parts representing the contributions from the major sources of
variation. The procedure uses these parts to test hypotheses concerning the equality of
mean effects of one or more factors and interactions between these factors. A special
use of the procedure enables the testing of the significance of regression, although in
simple linear regression this test may be calculated and presented in the form of a t-test.
Analysis of variance (ANOVA) is a process used to test the hypothesis that three or more
different samples were all selected from populations with the same mean. The method
is based on a test statistic:
𝑛𝑆 ∗ 2
𝐹= 2
𝑆
where n is the number of members in each sample, 𝑆 ∗ 2 is the variance of the sample
averages for all of the groups, and 𝑆 2 is the average variance for the groups. If the null
hypothesis is true and the population means actually are all the same, this statistic will
have an 𝐹 distribution with (𝑚 − 1) and 𝑚(𝑛 − 1) degrees of freedom, where 𝑚 is the
number of samples. If the value of the test statistic is too large, the null hypothesis is
rejected. Intuitively, a large value of 𝑆 ∗ 2 means that the observed sample averages are
spread further apart, thereby making the test statistic larger and the null hypothesis
less likely to be accepted.
ANALYTIC FUNCTION (COMPLEX ANALYSIS): A single-valued function which is
differentiable at every point of its domain. If the function is not differentiable
everywhere it may be said to be analytic at points where it is differentiable.
ANALYTIC GEOMETRY: Analytic geometry is the branch of mathematics that uses
algebra to help in the study of geometry.
ANALYTIC METHOD: If the OR, model is solved by using all the tools of classical
mathematics such as differential calculus and finite differences available for this task,
then such type of solutions are called analytic solutions. Solutions of various inventory
models are obtained by adopting the so called analytic procedure.

ANALYTIC SETS: We say that a subset 𝐴 of a complex analytic manifold 𝑮 is an analytic

set in 𝐺 if it is a closed subset and each point of 𝐴 has a neighborhood 𝑈 such that 𝑈 ∩ 𝐴
is the set of common zeros of a finite number of holomorphic functions in 𝑈. Specifically,
if 𝐴 is locally the set of zeros of a single holomorphic function that does not vanish
identically, then 𝐴 is called the principal analytic set.
ANCHOR RING: A surface obtained by rotating a circle (radius a) about a line in its plane
at a distance b (>a) from the centre is called anchor ring.

AND: The word “AND” is a connective word used in logic. The sentence “ p AND q” is true
only if both sentence p as well as sentence q are true. The operation of AND is illustrated
by the truth table:
p Q p AND q
T T T
T F F
F T F
F F F

An AND sentence is also called a conjunction.

ANGLE BETWEEN A LINE AND A PLANE: The complement of the angle between the
normal to a plane and a straight line is called the angle between the line and the plane.

Let the given plane be 𝑟 ∙ 𝑛 = 𝑞 and the given line ne parallel to the vector b. Then if ∅
be the angle between the normal to the plane i.e. n and the given line (parallel to b), we
have

𝑛 ∙ 𝑏 = 𝑛 𝑏 cos ∅ = 𝑛𝑏 cos ∅ , 𝑤𝑕𝑒𝑟𝑒 𝑛 = 𝑛 𝑎𝑛𝑑 𝑏 = 𝑏 .

𝑛 ∙𝑏
∴ ∅ = 𝑐𝑜𝑠 −1 .
𝑛𝑏

ANGLE BETWEEN TWO LINES: If a line has slope m, then the angle it makes with the
positive direction of x-axis is tan−1 𝑚.
The angle between a line with slope 𝑚1 and another line with slope 𝑚2 is given by
𝑚1 − 𝑚2
tan−1 .
1 + 𝑚1 𝑚2
If 𝒗𝟏 is a vector pointing in the direction of a line, and 𝒗𝟐 is a vector pointing in the
direction of another line, then the angle between them is:
𝒗𝟏 . 𝒗𝟐
cos−1
𝒗𝟏 × 𝒗𝟐
ANGLE BETWEEN TWO PLANES: The angle between the normals to two planes is called
the angle between the planes.

Let the two planes be 𝑟 ∙ 𝑛1 = 𝑞1 𝑎𝑛𝑑 𝑟 ∙ 𝑛2 = 𝑞2 𝑤𝑕𝑒𝑟𝑒 𝑛1 𝑎𝑛𝑑 𝑛2 are vectors

perpendicular to the two planes. Now if 𝜃 be the angle between the two planes (i.e.
between their normal’s, then 𝑛1 ∙ 𝑛2 = 𝑛1 𝑛2 cos 𝜃

𝑛1∙ 𝑛2 𝑛1∙ 𝑛2
or cos 𝜃 = where 𝑛1 = 𝑛1 𝑎𝑛𝑑 𝑛2 𝑛2 or 𝜃 = 𝑐𝑜𝑠 −1 .
𝑛1𝑛2 𝑛1𝑛2

ANGLE BETWEEN LINES IN SPACE: If ⌌𝑙1 , 𝑚1 , 𝑛1 ⌍ and ⌌𝑙2 , 𝑚2 , 𝑛2 ⌍ are direction ratios for
directions along the lines, the angle θ between the lines is given by
𝑙1 𝑙2 + 𝑚1 𝑚2 + 𝑛1 𝑛2
𝜃 = cos −1
𝑙1 2 + 𝑚1 2 + 𝑛1 2 𝑙2 2 + 𝑚2 2 + 𝑛2 2

ANGLE BISECTOR THEOREM: A theorem which states that, given any triangle and any
one of its angle, the angle bisector divides its opposite side into
2 segments whose length are of the same ratio as the lengths of the
corresponding adjacent sides.

ANGULAR ACCELERATION: The time derivative of angular velocity, that is, the rate
(with respect of time) of increase of angular velocity.

ANGULAR FREQUENCY: Also known as angular speed. (So called because the measure is
essentially the same as how often - frequency - a certain unit of angle is achieved.). It is the
𝑑2𝑥
constant ω in the equation = −𝜔2 𝑥 for simple harmonic motion. The angular
𝑑𝑡 2

frequency ω is usually measured in radians per second. The frequency of the

oscillations is equal to ω/2π.

ANGULAR MEASURE: The angle between the line from the observer to object 1 and the
line from the observer to object 2. For example: Shortly after sunrise, you see that the
sun is only a few degrees above the horizon, while at noon it is elevated to nearly 90
degrees (or directly overhead). You are the observer, the sun is Object 1, and the
horizon is Object 2.
ANGULAR MOMENTUM: A measure of the state of a physical system that is related to its
spin or abiliy to spin. Analogous to linear momentum, angular momentum of a system is
conserved unless an external force is applied in a particular way. Suppose that the
particle P of mass m has position vector 𝒓 and is moving with velocity 𝒗 at any time t.
Then the angular momentum 𝑳 of P about the point A with position vector 𝒓’ is the
vector defined by 𝑳 = (𝒓 – 𝒓’) × 𝑚𝒗. It is the moment of the linear momentum about
the point A

ANGULAR SPEED: Angular speed is the magnitude of the angular velocity.

ANGULAR VELOCITY: Suppose that a particle P is moving in the plane, in a circle with
centre at the origin 𝑂 and radius 𝑟0 . Let (𝑟0 , 𝜃) be the polar coordinates of P. At an
elementary level, the angular velocity may be defined to be θ.
ANNULUS: It is the region between two concentric circles. If the circles have radii 𝑟 and
1
𝑟 + 𝑤, the area of the annulus is equal to 2𝜋𝑤(𝑟 + 𝑤). It is therefore the same as the
2

area of a rectangle of width 𝑤 and length equal to the circumference of the circle
midway in size between the two original circles.

ANNULUS OF CONVERGENCE: The set of convergence of a Laurent series is either an

open set of the form {𝑧 ∶ 0 ≤ 𝑟1 < |𝑧 − 𝑃| < 𝑟2 }, together with perhaps some or all of
the boundary points of the set, or a set of the form {𝑧 ∶ 0 ≤ 𝑟1 < |𝑧 − 𝑃| < +∞},
together with perhaps some or all of the boundary points of the set. Such an open set is
called an annulus centered at 𝑃.
ANTECEDENT: The "if" part of a conditional statement. It is the part of a conditional
statement such that if the antecedent is true, then the other part (the consequent) must
also be true. For example, the antecedent is the clause "he is good at math" in the
statement "If he is good at math, he must know what is."
ANTIDERIVATIVE: Given a real function 𝑓 in the domain 𝐷, any function 𝑔 such that
𝑔′(𝑥) = 𝑓(𝑥), for all 𝑥 ∇ 𝐷, is an antiderivative of 𝑓. If 𝑔1 and 𝑔2 are both
antiderivatives of a continuous function 𝑓, then 𝑔1 (𝑥) and 𝑔2 (𝑥) differ by a constant. In
that case, the notation ∫ 𝑓(𝑥)𝑑𝑥 may be used for an antiderivative of f, with the
understanding that an arbitrary constant can be added to any antiderivative.
ANTIDIAGONAL: For conventionally written matrices, the diagonal that runs from the
top right to thr bottom left. i.e. The "other" diagonal, as opposed to the main diagonal.
(The main diagonal is the diagonal that the elements "1" in an identity matrix runs
along.)

ANTIKYTHERA MECHANISM: The Antikythera mechanism is an analog computer that

was used by the Greeks more than 2000 years ago to locate and predict the positions of
celestial objects. When new, the device was made of bronze with a wooden frame and
was roughly the size of a wall clock. The name Antikythera comes from the island near
where the device was discovered in the remains of a shipwreck that occurred in or
around the year 76 B.C. Like a mechanical clock, the Antikythera mechanism has dials
along with a sophisticated system of gears. It bears inscriptions that have been
deciphered. Based on its construction and on the nature of the writing, scientists have
deduced that the Antikythera mechanism must have been used to measure astronomical
time based on cycles of the sun, the moon and the planets. In this way it could have been
employed to predict solar and lunar eclipses, tides and the recurrence of the seasons.
These functions would have made the device useful to farmers, seafarers and perhaps
even military commanders.

ANTIPODAL POINTS: Antipodal points are two points on a sphere that are at opposite
ends of a diameter.
ANTISYMMETRIC MATRIX: A matrix such that it's sum with its transpose matrix is
the zero matrix.

ANTISYMMETRIC RELATION: A binary relation ~ on a set S is called as antisymmetric if,

for all 𝑎, 𝑏 ∇ 𝑆, whenever 𝑎 ~ 𝑏 and 𝑏 ~ 𝑎, then 𝑎 = 𝑏.
ANTI-SYMMETRIC TENSOR: A tenser 𝐴𝑖𝑗 is said to be anti-symmetric (or skew-
symmetric) if 𝐴𝑖𝑗 = −𝐴𝑗𝑖 .

ANTITONE FUNCTION: A function 𝑓 from a partially ordered set 𝑆 to a partially

ordered set 𝑇 is antitone if 𝑥 ≤ 𝑦 in 𝑆 implies 𝑓(𝑦) ≤ 𝑓(𝑥) in 𝑇.
APOLLONIUS: Apollonius of Perga (262 BC to 190 BC) was a mathematician who
studied in Alexandria under pupils of Euclid, wrote works that extended Euclid’s work
in geometry, particularly focusing on conic sections.

APOTHEM: The apothem of a regular polygon is the distance from the center of the
polygon to one of the sides of the polygon, in the direction that is perpendicular to that
side.

APPROXIMATE EQUALITY: Approximate equality is a concept used primarily in physics

and engineering, and also occasionally in mathematics. Two quantities are
approximately equal when they are close enough in value so the difference is
inconsequential in practical terms. Approximate equality is symbolized by a squiggly
equal sign ( ).
As an example of how approximate equality can be used in mathematics, consider the
positive square root of 2 (or 2 ). This is an irrational number ; when written in
decimal form, it is nonterminating and nonrepeating. Expressed to four significant
digits, 2 1.414. It is possible to express it to many more significant digits, but the
result will always be an approximation. Thus, to 10 significant digits,
2 1.414213562, and to 20 significant digits, 2 1.4142135623730950488. An
example of an approximate equation using variables is x + y z.

In many cases, the ordinary equality symbol (=) is used in situations where, to be
rigorous, the approximate equality sign ought to be used. This is done for two reasons:
first, most fonts do not include an approximate equality symbol; and secondly, many
people do not know what the approximate equality symbol means.
APPROXIMATION: When two quantities 𝑋 and 𝑥 are approximately equal, written
𝑋 ≈ 𝑥, one of them may be used in suitable circumstances in place of, or as an
22
approximation for, the other. For example, 𝜋 = and 3 = 1.73.
7

APPROXIMATION SCHEME: An approximation method requires a set of approximating

functions, 𝑈 say, which is a subset of 𝑋. Specifically, the method is just a mapping from 𝑋
to 𝑈. In other words, given any 𝑓 ∇ 𝑋, the method picks the element 𝑢𝑓 , say, from 𝑈,
which is regarded as an approximation to 𝑓. To find whether the method is good or bad,
it should be compared with the best approximation
APPROXIMATION THEORY: It is the area of numerical analysis related to finding a
simpler function 𝑓(𝑥) which has approximately the same value as 𝑔(𝑥) over a specified
interval. The analysis is then carried out on the more accessible 𝑓(𝑥) knowing that the
difference will be small.
ARAKELOV THEORY: Arakelov theory is an approach to arithmetic geometry that
explicitly includes the 'infinite primes'.
ARBITRARY CONSTANT: It is a non–numerical symbol which is not a variable in a
generalized operation. For example, 𝑦 = 𝑚𝑥 + 𝑐 is the general equation of a straight
line in two dimensions, where 𝑚 and 𝑐 are arbitrary constants which represent the
slope of the line and the y- intercept.
ARC: An arc of a circle is the set of points on the circle that lie in the interior of a
particular central angle. Therefore an arc is a part of a circle. The degree measure of an
arc is the same as the degree measure of the angle that defines it. If 𝜃 is the degree
measure of an arc and r is the radius, then the length of the arc is 2𝜋𝑟𝜃/360. The term
arc is also used for a portion of any curve.
arccos: It is the abbreviation for inverse cosine function.
arccosec: It is the abbreviation for inverse cosecant function. (Also written as arccsc).
arccosech: It is the abbreviation for inverse hyperbolic cosecant function. (Also written
as arcsch).
arccosh: It is the abbreviation for inverse hyperbolic cosine function.
arccot: It is the abbreviation for It is the abbreviation for inverse cotangent function.
arccoth: It is the abbreviation for inverse hyperbolic cotangent function.
arccsc: It is the abbreviation for inverse cosecant function. (Also written as arccosec).
arccsch: It is the abbreviation for inverse hyperbolic cosecant function. (Also written as
arcosech).
ARC LENGTH: The length of an arc of a curve can be found with integration. Let ds
represent a very small segment of the arc, and let dx and dy represent the x and y
components of the arc.
Then: 𝑑𝑠 = 𝑑𝑥 2 + 𝑑𝑦 2 .

Rewrite this as:

2
𝑑𝑦
𝑑𝑠 = 1+ 𝑑𝑥
𝑑𝑥
Now, suppose we need to know the length of the arc between the lines 𝑥 = 𝑎 and 𝑥 =
𝑏. Set up this integral:
𝑏
2
𝑑𝑦
𝑑𝑠 = 1+ 𝑑𝑥
𝑑𝑥
𝑎

For the curve 𝑥 = 𝑥(𝑡), 𝑦 = 𝑦(𝑡)(𝑡 ∇ [𝑎, 𝑏]), the arc length equals
𝑡=𝑏
2 2
𝑑𝑥 𝑑𝑦
𝑑𝑠 = + 𝑑𝑡
𝑑𝑡 𝑑𝑡
𝑡=𝑎

For the curve 𝑟 = 𝑟(𝜃)(𝑎 ≤ 𝜃 ≤ 𝑏), the arc length equals

𝜃=𝑏
2
𝑑𝑟
𝑑𝑠 = 𝑟2 + 𝑑𝜃
𝑑𝜃
𝜃=𝑎

arcsec: It is the abbreviation for inverse secant function.

arcsin: It is the abbreviation for inverse sine function.
arctan: It is the abbreviation for inverse tangent function.
ARCHIMEDEAN ORDERED FIELD: An ordered field F is said to be an Archimedean
ordered field if ∀𝑥, 𝑦 ∇ 𝐹, 𝑦 > 0, there exists some 𝑛 ∇ 𝑁 such that 𝑛𝑦 > 𝑥.
ARCHIMEDEAN PROPERTY OF REAL NUMBERS: Let 𝑎 be any real number and 𝑏 any
positive real number. Then there exists a positive integer 𝑛 such that 𝑛𝑏 > 𝑎.
ARCHIMEDEAN SPIRAL: It is the curve whose equation is 𝑟 = 𝑎𝜃 in polar coordinates,
where 𝑎 (> 0) is a constant.

ARCHIMEDES: Archimedes (290 BC to 211 BC) studied at Alexandria and then lived in
Syracuse. He wrote extensively on mathematics and developed formulas for the volume
and surface area of a sphere, and a way to calculate the circumference of a circle. He also
developed the principle of floating bodies and invented military devices that delayed
the capture of the city by the Romans.
ARCHIMEDES’ SPIRAL INVERSE: Taking the origin as the inversion center, archimedes'
spiral 𝑇 = 𝑎𝜃 inverts to the hyperbolic spiral = 𝑎/𝜃 .

AREA: The size of a surface measured in square units is called area.

AREA OF A CIRCLE: The area of a circle with radius r is equal to r2.
AREA OF AN ELLIPSE: The area of an ellipse (semi-major axis a and semi-minor axis b)
is equal to ab.
AREA OF A PARALLELOGRAM: The area of a parallelogram is equal to the length of a
base multiplied by the height to that base.
AREA OF A SQUARE OR A RECTANGLE: The area of a rectangle equals the product of its
base and its height (or the product of its length and its width). The area of a square is
equal to the square of the length of a side.
AREA OF A SURFACE OF REVOLUTION: Let 𝑦 = 𝑓(𝑥) be the graph of a function 𝑓 such
that 𝑓′ is continuous on [𝑎, 𝑏] and 𝑓(𝑥) ≥ 0 for all 𝑥 in [𝑎, 𝑏]. The area of the surface
obtained by rotating, through one revolution about the 𝑥-axis, the arc of the curve
𝑦 = 𝑓(𝑥) between 𝑥 = 𝑎 and 𝑥 = 𝑏, equals
𝑏
2
𝑑𝑦
2𝜋 𝑦 1+ 𝑑𝑥
𝑑𝑥
𝑎

For the curve 𝑥 = 𝑥(𝑡), 𝑦 = 𝑦(𝑡)(𝑡 ∇ [𝑎, 𝑏]), the surface area equals
 𝑡=𝑏
2 2
𝑑𝑥 𝑑𝑦
2𝜋 𝑦 + 𝑑𝑡
𝑑𝑡 𝑑𝑡
𝑡=𝑎

For the curve 𝑟 = 𝑟(𝜃)(𝑎 ≤ 𝜃 ≤ 𝑏), the surface area equals

𝜃=𝑏
2
𝑑𝑟
2𝜋 𝑟 sin 𝜃 𝑟 2 + 𝑑𝜃
𝑑𝜃
𝜃=𝑎

AREA OF A TRAPEZOID: The area of a trapezoid is equal to half the product of its height
and the sum of its bases.
AREA OF A TRIANGLE: The area of a triangle is equal to half the product of a base (a
side) and the height to that base.
AREA-PRESERVING MAP: A map 𝐹 from 𝐼𝑘 n to 𝑅 n is area-preserving if 𝑚(𝐴) =
𝑚(𝐹 𝐴 ) for every subregion 𝐴 of 𝐼𝑘 n , where 𝑚(𝐴) is the measure of 𝐴. A linear
transformation is area preserving if its corresponding determinant is equal to 1.

AREA SAMPLING: A method of sampling where a geographical region is subdivided into

smaller areas (counties, villages, city blocks, etc.), some of which are selected at random,
and the chosen areas are then subsampled or completely surveyed.

AREA THEOREM (CONFORMAL MAPPING): Suppose that 𝑓 is analytic and injective in

the punctured open unit disk 𝐷 − {0} and has the power series representation

then the coefficients satisfy

AREA UNDER A CURVE: Suppose that the curve 𝑦 = 𝑓(𝑥) lies above the 𝑥 −axis, so that
𝑓(𝑥) ≥ 0 for all 𝑥 in [𝑎, 𝑏]. The area under the curve, that is, the area of the region
bounded by the curve, the 𝑥-axis and the lines 𝑥 = 𝑎 and 𝑥 = 𝑏, equals
𝑏

𝑓 𝑥 𝑑𝑥
𝑎
If 𝑓(𝑥) ≤ 0 for all 𝑥 in [𝑎, 𝑏], the integral above is negative. However, it is still the case
that its absolute value is equal to the area of the region bounded by the curve, the
𝑥 −axis and the lines 𝑥 = 𝑎 and 𝑥 = 𝑏. If 𝑦 = 𝑓(𝑥) crosses the 𝑥-axis, appropriate
results hold. For example, if the regions A and B are as shown in the figure below,

then area of region A is

𝑏

𝑓 𝑥 𝑑𝑥
𝑎

and area of region B is

𝑐

− 𝑓 𝑥 𝑑𝑥
𝑏

It follows that
𝑐 𝑏 𝑐

𝑓 𝑥 𝑑𝑥 = 𝑓 𝑥 𝑑𝑥 + 𝑓 𝑥 𝑑𝑥
𝑎 𝑎 𝑏

Similarly, to find the area of the region bounded by a suitable curve, the 𝑦-axis, and lines
𝑦 = 𝑐 and 𝑦 = 𝑑, an equation 𝑥 = 𝑔(𝑦) for the curve must be found. Then the
required area equals
 𝑑

𝑔 𝑦 𝑑𝑦
𝑐

If a curve has an equation 𝑟 = 𝑟(𝜃) in polar coordinates, there is an integral that gives
the area of the region bounded by an arc AB of the curve and the two radial lines OA and
OB. Suppose that ∠xOA = α and ∠xOB = β. The area of the region described equals
𝛽
1
𝑟 2 𝑑𝜃
2
𝛼

ARG: It is the abbreviation for argument of a complex number.

ARGAND, JEAN ROBERT (1768–1822): Jean Robert Argand was a Swiss Mathematician,
who invented a geometrical representation of complex numbers. The name Argand
diagram is due to him.
ARGAND DIAGRAM: Another name for the complex plane. In a Cartesian coordinate
system, a point can be represented using coordinates (𝑥, 𝑦). When this point is taken to
represent the complex number (𝑥 + 𝑖𝑦), the plane is called complex plane or Argand
diagram. It is named after the Swiss-born mathematician Jean Robert Argand, one of
several people who invented this geometrical representation for complex numbers.
ARG MAX: It is the abbreviation for argument of the maximum.
ARG MIN: It is the abbreviation for argument of the minimum.
ARGUMENT: (1) The argument of a function is the independent variable that is put into
the function. In the expression sin x, x is the argument of the sine function.
(2) In logic, an argument is a sequence of sentences (called premises) that lead to a
resulting sentence (called the conclusion).
ARISTOTLE: Aristotle (384 BC to 322 BC) wrote about many areas of human knowledge,
including the field of logic. He was a student of Plato and also a tutor to Alexander the
Great.
ARITHMETIC-GEOMETRIC MEAN INEQUALITY: The arithmetic mean of a set of non-
negative numbers is never less than their geometric mean. So for any 𝑎, 𝑏 > 0, we have
𝑎+𝑏
≥ 𝑎𝑏 with equality if and only 𝑎 = 𝑏.
2

ARITHMETIC MEAN: The arithmetic mean of a group of n numbers (𝑎1 , 𝑎2 , . . . 𝑎𝑛 ),

written as AM, is the sum of the numbers divided by n:
𝑎1 + 𝑎2 + . . . +𝑎𝑛
𝐴𝑀 =
𝑛
The arithmetic mean is commonly called the average.
ARITHMETIC SEQUENCE: An arithmetic sequence is a sequence of n numbers of the
form
𝑎, 𝑎 + 𝑏, 𝑎 + 2𝑏, 𝑎 + 3𝑏, . . . , 𝑎 + (𝑛 − 1)𝑏
ARITHMETIC SERIES: An arithmetic series is a sum of an arithmetic sequence:
𝑎 + 𝑎 + 𝑏 + 𝑎 + 2𝑏 + 𝑎 + 3𝑏 + . . . +( 𝑎 + 𝑛 − 1 𝑏)
In an arithmetic series the difference between any two successive terms is a constant
(in this case b). The sum of the first n terms in the arithmetic series above is
𝑛
𝑆𝑛 = 2𝑎 + 𝑛 − 1 𝑏
2
ARITY: The arity of something is the number of arguments it takes. This is usually
applied to functions: an 𝑛-ary function is one that takes 𝑛 arguments. Unary is a
synonym for 1-ary, and binary for 2-ary.
ARMSTRONG NUMBER: The 𝑛-digit numbers equal to sum of nth powers of their digits
(a finite sequence), also called plus perfect numbers. They first few are given by 1, 2, 3,
4, 5, 6, 7, 8, 9, 153, 370, 371, 407, 1634, 8208, 9474, 54748.

ARRAY: A configuration or arrangement of objects which is considered systematic in

some ways, including but not limited to a rectangular array, where matrices (and by
extension, the matrix representation of vectors) are examples.

ARTIFICIAL SINGULARITIES: These singularities are due to particular choice of

parametric representation e.g. the pole in the plane, referred to polar coordinates is an
artificial singularity.
ARTIN APPROXIMATION THEOREM: Let 𝐾 be a valued field of characteristic zero and let
𝑓(𝑥, 𝑦) be a vector of convergent power series in two sets of variables 𝑥 and 𝑦. Assume
given a formal power series solution 𝑦(𝑥) vanishing at 0, 𝑓(𝑥, 𝑦 (𝑥)) = 0. Then there
exists, for any 𝑐 ∇ 𝑁, a convergent power series solution 𝑦(𝑥), 𝑓(𝑥, 𝑦(𝑥)) = 0 which
coincides with 𝑦(𝑥) up to degree c, 𝑦(𝑥) ≡ 𝑦(𝑥) modulo (𝑥)𝑐 .
ARTIN, EMIL: Emil Artin (March 3, 1898-December 20, 1962) was born in Vienna. After
he studied at the Universities of Vienna, Leipzig, and Gottingen, he taught at the
University of Hamburg from 1923 to 1937. In 1937 he left Germany under the Nazi
regime for America, where he taught at Notre Dame, Indiana, and Princeton universities.
He returned to Germany in 1948 and taught again at the University of Hamburg until his
fatal heart attack.
ARTINIAN GROUP: A group in which any decreasing chain of distinct subgroups
terminates after a finite number.

ARTINIAN RING: A non-commutative semi-simple ring satisfying the “descending chain

condition.”

ARTIN L-FUNCTIONS: Artin L-functions are defined for quite general Galois
representations. The introduction of étale cohomology in the 1960s meant that Hasse–
Weil L-functions (q.v) could be regarded as Artin L-functions for the Galois
representations on l-adic cohomology groups.
ARTIN’S CONJECTURE: There are at least two statements which go by the name of
Artin’s conjecture. The first is the Riemann hypothesis. The second states that every
integer not equal to −1 or a square number is a primitive root modulo 𝑝 for infinitely
many p and proposes a density for the set of such 𝑝 which are always rational multiples
of a constant known as Artin's constant. There is an analogous theorem for functions
instead of numbers which has been proved by Billharz.

ARTIN’S RECIPROCITY THEOREM: A general reciprocity theorem for all orders. If R is

a number field and R’ a finite integral extension, then there is a surjection from the
group of fractional ideals prime to the discriminant, given by the Artin symbol. For some
cycle c, the kernel of this surjection contains each principal fractional ideal generated by
an element congruent to 1 mod c.

ARTIN–WEDDERBURN THEOREM: Artin–Wedderburn theorem is a classification

theorem for semi-simple rings and semi-simple algebras. The theorem states that an
(Artinian) semi-simple ring 𝑅 is isomorphic to a product of finitely many 𝑛𝑖 × 𝑛𝑖 matrix
rings over division rings 𝐷𝑖 , for some integers 𝑛𝑖 , both of which are uniquely determined
up to permutation of the index 𝑖. In particular, any simple left or right Artinian ring is
isomorphic to an 𝑛 × 𝑛 matrix ring over a division ring 𝐷, where both 𝑛 and 𝐷 are
uniquely determined.
ĀRYABHATA (about 476-550): An Indian mathematician, author of one of the oldest
Indian mathematical texts. Written in verse, the Āryabhatīya is a summary of
miscellaneous rules for calculation and mensuration. It deals with the areas of certain
plane figures, values for π, and the summation of arithmetic series.

ARZELA-ASCOLI THEOREM: A sequence {𝑓𝑛  } of continuous functions on an interval 𝐼 =

[𝑎, 𝑏] is uniformly bounded if there is a number 𝑀 such that

for every function  𝑓𝑛   belonging to the sequence, and every 𝑥 ∇ [𝑎, 𝑏]. The sequence
is equicontinuous if, for every 𝜀 > 0, there exists 𝛿 > 0 such that

whenever |𝑥 − 𝑦| < 𝛿  for all functions  𝑓𝑛   in the sequence. Succinctly, a sequence is
equicontinuous if and only if all of its elements admit the same modulus of continuity. In
simplest terms, the theorem can be stated as follows:

Consider a sequence of real-valued continuous functions  𝑓𝑛   ; 𝑛 ∇ 𝑵 defined on a

closed and bounded interval [𝑎, 𝑏] of the real line. If this sequence is uniformly
bounded and equicontinuous, then there exists a subsequence (𝑓𝑛𝑘 ) that converges
uniformly.
The converse is also true, in the sense that if every subsequence of  𝑓𝑛  } itself has a
uniformly convergent subsequence, then  𝑓𝑛  } is uniformly bounded and
equicontinuous.
ASCENDING CHAIN CONDITION: A collection 𝑆 of subsets of a set 𝑋 satisfies the
ascending chain condition or ACC if there does not exist an infinite ascending chain
𝑆1 ⊆ 𝑆2 ⊆ · · · of subsets from 𝑆.

ASSIGNMENT PROBLEM: It is a type of problem in which things of one type are to be

matched with the same number of things of another type in a way that is, in a specified
sense, the best possible.
 The number of assignees and the number of tasks are the same, 𝑛.
  Each assignee is to be assigned to exactly one task.
 Each task is to be performed by exactly one assignee.
 There is a cost 𝑐𝑖𝑗 associated with assignee 𝑖 (𝑖 = 1,2, … , 𝑛) performing
task 𝑗 (𝑗 = 1,2, … , 𝑛).
 The objective is to determine how all 𝑛 assignments should be made to minimize
the total cost.
ASSOCIATED LAGUERRE POLYNOMIALS: The differential equation

d2 y dy
𝑥 + 𝛼 + 1 − 𝑥 + ny = 0,
dx 2 dx

is called associated Laguerre equation.

ASSOCIATION: A relationship between two or more variables.

ASSOCIATIVE: Denoting an operation in which the outcome is independent of the
grouping of the symbols and numbers involved, that is,
(𝑎 ∗ 𝑏) ∗ 𝑐 = 𝑎 ∗ (𝑏 ∗ 𝑐), in which ∗ is the operation. Such examples include addition
and multiplication.

(1 + 2) + 3 = 1 + (2 + 3)
(4 × 5) × 6 = 4 × (5 × 6)

The following examples demonstrate that subtraction and division are not associative:

(3 − 2) − 1 3 − (2 − 1)
(12 / 4) / 3 12 / (4 / 3)

ASSOCIATIVE ALGEBRAS: Let 𝐾 be a commutative ring with unity element 1, and let 𝐴
be a ring which is a unitary 𝐾-module. Such a ring 𝐴 is called an associative algebra over
𝐾 (or simply algebra over 𝐾 ) if it satisfies the condition 𝛼(𝑎𝑏) = (𝛼𝑎)𝑏 = 𝑎(𝛼𝑏) (𝛼 ∇
𝐾; 𝑎, 𝑏 ∇ 𝐴). An (associative) algebra 𝐴 over 𝐾 is often written 𝐴/𝐾, and 𝐾 is called the
coefficient ring (or ground ring) of the algebra 𝐴 = 𝐴/𝐾 . In particular, if 𝐾 is a field,
then it is called the coefficient field (or ground field) of 𝐴. Notions such as zero algebra,
unitary algebra, commutative algebra, (semi) simple algebra, and division algebra are
replicas of the respective ones for rings. Considering both structures as rings and as 𝐾 -
modules, homomorphisms and isomorphisms are defined in a natural manner, and are
called algebra homomorphism and algebra isomorphism, respectively. In this
connection, subalgebra, quotient algebra (or residue class algebra), and direct product
of algebras are also defined as in the case of rings.
ASSOCIATIVE BINARY OPERATIONS: Let ∗ be a binary operation on a set 𝐴. Given any
three elements 𝑥, 𝑦 and 𝑧 of a set 𝐴, the binary operation, applied to the elements 𝑥 ∗ 𝑦
and 𝑧 of 𝐴, yields an element (𝑥 ∗ 𝑦) ∗ 𝑧 of 𝐴, and, applied to the elements 𝑥 and 𝑦 ∗ 𝑧
of 𝐴, yields an element 𝑥 ∗ (𝑦 ∗ 𝑧) of 𝐴. A binary operation ∗ on a set 𝐴 is said to be
associative if (𝑥 ∗ 𝑦) ∗ 𝑧 = 𝑥 ∗ (𝑦 ∗ 𝑧) for all elements 𝑥, 𝑦 and 𝑧 of 𝐴. The
operations of addition and multiplication on the set 𝑅 of real numbers are associative,
since (𝑥 + 𝑦) + 𝑧 = 𝑥 + (𝑦 + 𝑧) and (𝑥 × 𝑦) × 𝑧 = 𝑥 × (𝑦 × 𝑧) for all real numbers
𝑥, 𝑦 and 𝑧. However the operation of subtraction is not associative. For example
(1 − 2) − 3 = −4, but 1 − (2 − 3) = 2.
ASSOCIATIVE PROPERTY: An operation obeys the associative property if the grouping of
the numbers involved does not matter. Formally, the associative property of addition
says that
𝑎+ 𝑏 + 𝑐= 𝑎+ 𝑏+ 𝑐 .
The associative property for multiplication says that
(𝑎 × 𝑏) × 𝑐 = 𝑎 × 𝑏 × 𝑐 .
ASSUMPTION OF SEQUENCING: Some principal assumptions are as follows:

(i) The processing times 𝐴′𝑖 𝑠 etc. are exactly known and are independent of the
order of the jobs in which they are to be processed. Such problems where times are
exactly known are called deterministic problems.
(ii) The time taken by the jobs in going from one machine to another is negligible.
(iii) Each job, once started on a machine, is to be performed up to completion on that
machine.
(iv) A job starts on the machine as soon as the job and the machine both are idle and
job is next to the machine and the machine is also next to the job.
(v) There is only one machine of each type.
(vi) No machine may process more than one job at a time.
(vii) The cost of keeping the jobs in inventory (if needed) during the in process is
same for all jobs. Also it is too small that it can be neglected.
(viii) The order of completion of jobs has no significance i.e., no job is to be given
priority. Times of jobs are independent of sequence of jobs.

ASTROID: It is a hypocycloid in which the radius of the rolling circle is a quarter of the
radius of the fixed circle. It has parametric equations 𝑥 = 𝑎 𝑐𝑜𝑠 3 𝑡, 𝑦 = 𝑎 𝑠𝑖𝑛3 𝑡, where
𝑎 is the radius of the fixed circle.

ASYMMETRIC RELATION: A binary relation ~ on a set S is said to be asymmetric if,

∀ 𝑎, 𝑏 ∇ 𝑆, whenever 𝑎 ~ 𝑏, then 𝑏 ~ 𝑎 does not hold. For example, the relation > on
the set of integers is asymmetric.
ASYMPTOTE: An asymptote is a straight line that is a close approximation to a particular
curve as the curve goes off to infinity in one direction. The curve becomes very, very
close to the asymptote line, but never touches it.

ASYMPTOTIC: 1. Of or related to an asymptote, having an asymptote or capable of

having an asymptote.

2. The relation between two functions that tends to the same value (possibly infinite)
and such that the difference between the values of the 2 functions, considered as a
proportion of the values of either function, becomes arbitrarily small.

ATKINSON'S THEOREM: Let 𝐻 be a Hilbert space and 𝐿(𝐻) the set of bounded operators
on 𝐻. An operator 𝑇 ∇ 𝐿(𝐻) is said to be a Fredholm operator if the kernel 𝐾𝑒𝑟(𝑇) is
finite-dimensional, 𝐾𝑒𝑟(𝑇 ∗ ) is finite-dimensional (where 𝑇 ∗ denotes the adjoint of 𝑇),
and the range 𝑅𝑎𝑛(𝑇) is closed. Atkinson's theorem states:
A 𝑇 ∇ 𝐿(𝐻) is a Fredholm operator if and only if 𝑇 is invertible modulo compact
perturbation, i.e. 𝑇𝑆 = 𝐼 + 𝐶1 and 𝑆𝑇 = 𝐼 + 𝐶2 for some bounded
operator 𝑆 and compact operators 𝐶1 and 𝐶2 .
ATLASES: The description of most manifolds requires more than one chart (a single
chart is adequate for only the simplest manifolds). A specific collection of charts which
covers a manifold is called an atlas. An atlas is not unique as all manifolds can be
covered multiple ways using different combinations of charts. Two atlases are said to
be Ck equivalent if their union is also a Ck atlas.
AUGEND: An argument in the binary operation of addition. Technically the second
argument, it is commonly applied to both arguments due to addition
being commutative.

AUGMENTED MATRIX: For a given set of 𝑚 linear equations in 𝑛 unknowns 𝑥1 , 𝑥2 , … 𝑥𝑛 ,

the augmented matrix is the matrix obtained by adjoining to the matrix of coefficients
an extra column of entries taken from the right-hand sides of the equations.

We would like to make the process of solving a system of linear equations more
mechanical by forgetting about the variable names 𝑤, 𝑥, 𝑦, 𝑧 etc and doing the whole
operation as a matrix calculation. For this, we use the augmented matrix of the system
of equations, which we have constructed by “glueing” an extra column on the right-hand
side of the matrix representing the linear transformation, as above.

For the system 𝐴𝑥 = 𝐵 of m equations in n unknowns, where A is the 𝑚 × 𝑛 matrix (𝑎𝑖𝑗 )

is defined to be the 𝑚 × (𝑛 + 1) matrix.
 𝛼11 𝛼12 ⋯ 𝛼1𝑛 : 𝑏1
𝛼21 𝛼22 ⋯ 𝛼2𝑛 : 𝑏2
⋮ ⋮ ⋱ ⋮ : ⋮
𝛼𝑚1 𝛼𝑚2 ⋯ 𝛼𝑚𝑛 : 𝑏𝑚

The vertical line in the matrix is put there to remind that the rightmost column is
different from the others, and arises from the constants on the right hand side of the
equations.

Let us look at the following system of linear equations over ℝ: Suppose that we want to
find all 𝑤, 𝑥, 𝑦, 𝑧 𝜖 ℝ satisfying the equations.
2𝑤 − 𝑥 + 4𝑦 − 𝑧 = 1
𝑤 + 2𝑥 + 𝑦 + 𝑧 = 2
𝑤 − 3𝑥 + 3𝑦 − 2𝑧 = −1
−3𝑤 − 𝑥 − 5𝑦 + 0𝑧 = −3

Elementary row operations on 𝐴 are precisely Gauss transformations of the

corresponding linear system. Thus the solution can be carried out mechanically as
follows:

Matrix Operation(s)

2 −1 4 −1 : 1
1 2 1 1 : 2 𝑅1 →
𝑅1
1 −3 3 −2 : −1 2
−3 −1 −5 0 : −3

1 −1 2 2 −1 2 : 1 2 𝑅2 → 𝑅2 − 𝑅1
1 2 1 1 : 2 𝑅3 → 𝑅3 − 𝑅1
1 −3 3 −2 : −1
𝑅4 → 𝑅4 + 3𝑅1
−3 −1 −5 0 : −3

1 −1 2 2 −1 2 : 1 2
0 5 2 −1 3 2 : 3 2 𝑅3 → 𝑅3 + 𝑅2
0 −5 2 1 −3 2 : −3 2 𝑅4 → 𝑅4 + 𝑅2
0 −5 2 1 −3 2 : −3 2

1 −1 2 2 −1 2 : 1 2
0 5 2 −1 3 2 : 3 2 𝑅2 →
2𝑅2
0 0 0 0 : 0 5
0 0 0 0 : 0
 1 −1 2 2 −1 2 : 1 2
0 1 −2 5 3 5 : 3 5 𝑅1 → 𝑅1 +
𝑅2
0 0 0 0 : 0 2
0 0 0 0 : 0

1 0 9 5 −1 5 : 4 5
0 1 −2 5 3 5 : 3 5
0 0 0 0 : 0
0 0 0 0 : 0

The original system has been transformed to the following equivalent system, i.e, Both
systems have the same solutions.

𝑤 + 9𝑦 5 − 𝑧 5 = 4 5 , 𝑥 − 2𝑦 5 + 3 𝑧 5 = 3 5

In a solution to the latter, variables y and z can take arbitrary values in ℝ;

Say 𝑦 = 𝛼, 𝑧 = 𝛽.Then the equations tell us that 𝑤 = −9𝛼 5 + 𝛽 5 + 4 5 and

𝑥 = 2𝛼 5 − 3 𝛽 5 + 3 5 and so the complete set of solutions is

𝑤, 𝑥, 𝑦, 𝑧 = −9𝛼 5 + 𝛽 5 + 4 5 , 2𝛼 5 − 3 𝛽 5 + 3 5 , 𝛼, 𝛽

= (4 5 , 3 5,0,0) + 𝛼(−9 5, 2 5,1,0) + 𝛽(1 5, −3 5,0,1)

AUT: It is the abbreviation for automorphism group.

AUTOMORPHIC FUNCTION: A function 𝑓(𝑧) is said to be automorphic under a group of
transformations if it is analytic in a domain D except at singular points and for all
transformations T in the group that if z is in D then so is T(z).
AUTOMORPHIC NUMBER: A number 𝑘 such that 𝑛𝑘 2 has its last digits equal to 𝑘 is
called n-automorphic. For example, 1. 52 = 25 - and 1. 62 = 36 are 1-automorphic and
2 . 82 = 128 and 2 . 882 = 15488 are 2-automorphic. de Guerre and Fairbairn (1968)
gave a history of automorphic numbers. The first few l-automorphic numbers are 1, 5, 6,
25, 76, 376, 625, 9376, 90625, . . .

AUTOMORPHISM: An automorphism is a one-to-one correspondence mapping the

elements of a set onto itself, so the domain and range of the function are the same. For
example, 𝑓(𝑥) = 𝑥 + 5 is an automorphism on R but 𝑔(𝑥) = 𝑐𝑜𝑠 𝑥 is not.
1 1
AUXILLIARY SERIES: The series is known as the auxiliary series. Note that
𝑛𝑝 𝑛𝑝

converges if 𝑝 > 1 and diverges if 𝑝 < 1.

AVERAGE: The average of a group of numbers is the same as the arithmetic mean.
AVERAGE EXPECTED PAYOFF: An estimate of the amount that will be gained in a game
of chance, calculated by multiplying the probability of winning by the number of points
won each time.
AVERAGE LENGH OF LINE: The average time for which the system remains idle.

AXIAL PLANE: A fixed reference plane containing two of the three coordinate axes in a
three-dimensional coordinate system. For example, the plane that contains 𝑥-axis and 𝑦-
axis is the axial plane called the 𝑥𝑦-plane, or (𝑥, 𝑦)-plane. Similarly, the plane containing
𝑦-axis and 𝑧-axis is the 𝑦𝑧-plane.
AXIOM: An axiom is a statement that is assumed to be true without proof. Axiom is a
synonym for postulate.
AXIOMATIC SET THEORY: Axiomatic set theory pursues the goal of reestablishing the
essentials of G. Cantor’s rather intuitive set theory by axiomatic constructions
consistent with modern theories of the foundations of mathematics. A system of axioms
for set theory was first given by E. Zermelo, and was completed by A. Fraenkel. J. von
Neumann expressed it in symbolic logic, gave a formal generalization, and eliminated
ambiguous concepts. P. Bernays and K. Godel further refined and simplified von
Neumann’s formulation. The theories based on the systems before and after the formal
generalization are called Zermelo-Fraenkel set theory (ZF) and Bernays-Gode1 set
theory (BG), respectively.
AXIOM OF CHOICE: It states that for any collection of mutually exclusive sets there is at
least one set that contains exactly one element in common with each of the non-empty
sets.

∀𝑥 ∇ 𝑎∃𝑦𝐴(𝑥, 𝑦) → ∃𝑦∀𝑥 ∇ 𝑎𝐴(𝑥, 𝑦 𝑥 ).

This asserts that if for any element 𝑥 of a there is a set 𝑦 such that 𝐴(𝑥, 𝑦), then there is a
choice function 𝑦 for the formula, i.e., 𝐴 𝑥, 𝑦 𝑥 for all 𝑥 in a. Usually a function is
represented by its graph. A set 𝑓 is called a function defined on a if

∀𝑥∀𝑦(𝑥, 𝑦) ∇ 𝑓 → 𝑥 ∇ 𝑎), ∀𝑥 ∇ 𝑎∃𝑦(𝑥, 𝑦) ∇ 𝑓).

∀𝑥 ∇ 𝑎∀𝑦𝐴𝑧 𝑥, 𝑦 ∇ 𝑓⋀ 𝑥, 𝑧 ∇ 𝑓 → 𝑦 = 𝑧).
This formula is denoted by Fnc(𝑓); then the formula 𝐴 𝑥, 𝑓 𝑥 is an abbreviation of
Fun 𝑓 ⋀ 3𝑦(𝑥, 𝑦) ∇ 𝑓 ⋀ 𝐴 (𝑥, 𝑦)).

AXIOM OF THE EMPTY SET: ∃𝑥∀𝑦 ℸ𝑦 𝜖 𝑥 .

This asserts the existence of the empty set. The empty set is denoted by ⊘ or 0.

AXIOM OF EXTENSIONALITY: ∀𝑥 𝑥 𝜖 𝑎 ≡ 𝑥 𝜖 𝑏 → 𝑎 = 𝑏.

This asserts that sets formed by the same elements are equal. The formula 𝑥 𝜖 𝑎 (𝑥 𝜖 𝑏)
is denoted by 𝑎 ⊂ 𝑏. This means “𝑎 is a subset of 𝑏.” Then Axiom 1 can be expressed by

𝑎 ⊂ 𝑏 ∧ 𝑏 ⊂ 𝑎 → 𝑎 = 𝑏.

AXIOM OF INFINITY: ∃𝑥 0 𝜖 𝑥 ∨ ∀ 𝑦 𝜖 𝑥(𝑦 ′ 𝜖𝑥) .

This asserts the existence of the set consisting all the natural numbers, whereas
0,1 = 0′ = 0 , 2 = 1′ = 0,1 , 3 = 2′ = 0,1,2 . This definition of natural number is due
to von Neumann.

AXIOM OF THE POWER SET: ∀𝑥∃𝑦 𝑦 𝜖 𝑥 ≡ ∀𝑧 𝜖 𝑦(𝑧𝜖𝑎) .

This asserts the existence, for any sets a of the power set 𝑥 consisting of all the subsets
of a. This 𝑥 is denoted by 𝑃𝑎. We have 𝑆(𝑃 𝑎 = 𝑎, 𝑠𝑜 𝑆 is a left inverse of 𝑃 and the
products 𝑆𝑃 𝑎𝑛𝑑 𝑃𝑆 are idempotent.

AXIOM OF REGULARITY: ∀𝑥(∀𝑦𝜖𝑥𝐴(𝑦) → ∀(𝑥) → ∀ 𝑥𝐴 (𝑥).

Using this we can show that ℸ𝑥𝜖𝑥,

ℸ(𝑥𝜖𝑦 𝐴 𝑦𝜖 𝑥), etc. If we assume the axiom of choice stated below, then this is equivalent
to the nonexistence of an infinite descending sequence

𝑥𝑛 ∇ ⋯ ∇ 𝑥2 ∇ 𝑥1 .

If a model e1 of a set theory satisfies the axiom of regularity and has an infinite
descending sequence that is not in the model, then such a model is called a nonstandard
model.

AXIOM OF REPLACEMENT: ∃𝑥∀𝑦𝜖𝑎 ∃𝑧 𝐴(𝑦, 𝑧) → ∃𝑧 𝜖 𝑥𝐴 (𝑦, 𝑧) .

This asserts the existence for any set a of a set 𝑥 such that for any 𝑦 of 𝑎𝑖 , if there exists a
𝑧 satisfying 𝐴(𝑦, 𝑧) then such 𝑧 exists in 𝑥. If the relation 𝐴(𝑦, 𝑧) determines a function,
then the image of a set by the relation is included in a set, so by Axiom 7, the image is a
set. This axiom was introduced by Frankel.

AXIOM OF SEPARATION: ∃𝑥∀𝑦 𝑦 𝜖 𝑥 ≡ 𝑦 𝜖 𝑎 ⋀ 𝐴 (𝑦) .

This asserts that the existence for any set a and a formula 𝐴(𝑦) of a set 𝑥 consisting of all
element of a satisfying 𝐴(𝑦). This 𝑥 is denoted by {𝑦𝜖𝑎 𝐴 𝑦 }. This is rather a schema for
an infinite number of axioms, for there are an infinite number of 𝐴(𝑦). This axiom, also
called the axiom of comprehension or axiom of subsets, was introduced by Zermelo.

For example, the set of all natural numbers is introduced by

{𝑥𝜖𝑎}∀𝑦(0𝜖𝑦 ⋀ 𝐴 𝑧𝜖𝑦 (𝑧 ′ 𝜖𝑦) → 𝑥𝜖𝑦)}.

AXIOM OF THE SUM SET: ∃𝑥∀𝑦 𝑦 𝜖 𝑥 ≡ ∃𝑧 𝜖 𝑎(𝑦𝜖𝑧) .

This asserts the existence, for any sets a of the sum (or union) 𝑥 of all the sets that are
elements of a. This 𝑥 is denoted by the 𝑈𝑎 or 𝑆(𝑎). We write a 𝑈𝑏 for 𝑈 𝑎, 𝑏 and
𝑎′ for 𝑎𝑈{𝑎. 𝑎}.

AXIOM OF THE UNORDERED PAIR: ∃𝑥∀𝑦 𝑦 𝜖 𝑥 ≡ 𝑦 𝜖 𝑎 ∨ 𝑦 = 𝑏 .

This asserts the existence, for any sets a and b, of a set 𝑥 having a and b as its only
elements. This 𝑥 is called the unordered pair of a and b and is denoted by 𝑎, 𝑏 . The
notion of ordered pair is characterized by

𝑎, 𝑏 = 𝑐. 𝑑 ≡ 𝑎 = 𝑐 ∨ 𝑏 = 𝑑.

An example of such is 𝑎, 𝑏 = 𝑎, 𝑎 , 𝑎, 𝑏 .

AXIOM RELATING ADDITION AND MULTIPLICATION:

α + β . γ = α. γ + β. γ for all α, β, γ ϵ S (Distributivity)

AXIOMS OF PROBABILITY: There are three axioms of probability:

1) The probability of an event is a real number greater than or equal to zero; 2) The sum
of the probabilities of all possible outcomes of a given experiment is 1;
3) If two events cannot both occur at the same time, the chance that either one occurs is
the sum of the chances that each occurs.
AXIOMS FOR NUMBER SYSTEM: Any particular axiom might be true in some number
systems but not in others.

AXIOMS FOR ADDITION: Let S be any number system.

A1.α + β = β + α for all α, β ϵ S (Commutativity)

A2. α + β + γ = α + β + γ for all α, β, γ ϵ S (Associativity)

A3.There is a number 0 ∇ 𝑆 𝑠𝑢𝑐𝑕 𝑡𝑕𝑎𝑡 𝛼 + 0 = 0 + 𝛼 𝑓𝑜𝑟 𝑎𝑙𝑙 𝛼 𝜖 𝑆

(Identity)

A4.For each number 𝛼 𝜖 𝑆,there exist a number – 𝛼 𝜖 𝑆 𝑠𝑢𝑐𝑕 𝑡𝑕𝑎𝑡

𝛼 + −𝛼 = 0 = −𝛼 + 𝛼 (Inverse)

These axioms may or may not be satisfied by a given number system S.For example,
in ℕ, A1 and A2 hold but A3 and A4 do not hold. A1-A4 all hold in ℞, ℚ, ℝ 𝑎𝑛𝑑 ℂ.

AXIOMS FOR MULTIPLICATION:

M1.α. β = β. α for all α, β ϵ S (Commutativity)

M2. α. β .γ = α. β. γ for all α, β, γ ϵ S (Associativity)

M3. There is a number 1𝜖 𝑆 𝑠𝑢𝑐𝑕 𝑡𝑕𝑎𝑡 𝛼. 1 = 1. 𝛼 = 𝛼 for all α ϵ S (Identity)

M4. For each number 𝛼 𝜖 𝑆 𝑤𝑖𝑡𝑕 𝛼 ≠ 0,there exist a number 𝛼 −1 𝜖 𝑆 𝑠𝑢𝑐𝑕 𝑡𝑕𝑎𝑡

𝛼. 𝛼 −1 = 1 = 𝛼 −1 . 𝛼 (Inverse)

In ℕ 𝑎𝑛𝑑 ℞, M1-M3 hold but M4 does not hold. M1-M4 all hold in ℚ, ℝ 𝑎𝑛𝑑 ℂ.

AXIS OF SYMMETRY: It is a line in which the two halves of a curve reflect into each
other.
AXIS-SYMMETRIC FLOW FIELD: A flow field is said to be axis-symmetric when the
velocity components (𝑢, 𝑣, 𝑤) with regard to cylindrical coordinates x, r, ∅ are all
independent of the azimuthal angle ∅. The velocity component q ∅ in the direction of
increasing ∅ is called swirl component of velocity.

AZIMUTH: The angle between the projection of the radius vector onto the plane
perpendicular to the polar axis, and the initial line.

B
BABBAGE, CHARLES (1792–1871): Charles Babbage was a British mathematician and
inventor of mechanical calculators.
BACK-SUBSTITUTION: A method of solving a system of linear equations (or
the equivalent matrix equation) that involves finding the components in the reverse
order to the way it is indexed, due to its being manipulated to a row-echelon form. (The
Upper-triangular form of a square matrix is a special case of this.)

BACKWARD DIFFERENCE: If {(𝑥𝑖 , 𝑓𝑖 )}, 𝑖 = 0, 1, 2, … is a given set of function values with

𝑥𝑖+1 = 𝑥𝑖 + 𝑕, 𝑓𝑖 = 𝑓(𝑥𝑖 ) then the backward difference at 𝑓𝑖 is defined by
𝑓𝑖 – 𝑓𝑖−1 = 𝑓 𝑥𝑖 − 𝑓(𝑥𝑖−1 ).

BAIRE CATEGORY THEOREM: A Baire space is a topological space with the following
property: for each countable collection of open dense sets 𝑈𝑛 , their intersection ⋂ 𝑈𝑛 is
dense.

(BCT1) Every complete metric space is a Baire space. More generally, every topological
space which is homeomorphic to an open subset of a complete pseudometric space is a
Baire space. Thus every completely metrizable topological space is a Baire space.

(BCT2) Every locally compact Hausdorff space is a Baire space. The proof is similar to
the preceding statement; the finite intersection property takes the role played by
completeness.

Note that neither of these statements implies the other, since there are complete metric
spaces which are not locally compact (the irrational numbers with the metric defined
below; also, any Banach space of infinite dimension), and there are locally compact
Hausdorff spaces which are not metrizable (for instance, any uncountable product of
non-trivial compact Hausdorff spaces is such; also, several function spaces used in
Functional Analysis; the uncountable Fort space).

(BCT3) A non-empty complete metric space is not the countable union of nowhere-
dense closed sets.

This formulation is equivalent to BCT1 and is sometimes more useful in applications.

Also: if a non-empty complete metric space is the countable union of closed sets, then
one of these closed sets has non-empty interior.

BAIRE SPACE: This has two distinct meanings of Baire spaces:

 A space is a Baire space if the intersection of any countable collection of dense

open sets is dense;
 Baire space is the set of all functions from the natural numbers to the natural
numbers, with the topology of pointwise convergence
BALINSKI'S THEOREM: It is a theorem about the graph-theoretic structure of three-
dimensional polyhedra and higher-dimensional polytopes. It states that, if one forms
an undirected graph from the vertices and edges of a convex d-dimensional polyhedron
or polytope (its skeleton), then the resulting graph is at least d-vertex-connected: the
removal of any d − 1 vertices leaves a connected subgraph. For instance, for a three-
dimensional polyhedron, even if two of its vertices (together with their incident edges)
are removed, for any pair of vertices there still exists a path of vertices and edges
connecting the pair.
BALKING: A customer may not like to join the queue seeing it very long and he may not
like to wait.

BALL: A ball in a metric space is a set containing all the points which are not more than a
given distance (radius of the ball) from a fixed point (center of the ball). An open ball
does not include the boundary where the distance is equal to the given constant and a
closed ball does include the boundary.
BANACH–ALAOGLU THEOREM : Banach–Alaoglu theorem also known as Alaoglu's
theorem states that the closed unit ball of the dual space of a normed vector
space is compact in the weak topology.
BANACH FIXED POINT THEOREM: Let (𝑋, 𝑑) be a non-empty complete metric
space with a contraction mapping 𝑇 : 𝑋 → 𝑋. Then 𝑇 admits a unique fixed-
point 𝑥 ∗ in 𝑋 (i.e. 𝑇(𝑥 ∗ ) = 𝑥 ∗ ). Furthermore, 𝑥 ∗ can be found as follows:
Start with an arbitrary element 𝑥0 in X and define a sequence {𝑥𝑛 } by 𝑥𝑛 = 𝑇(𝑥𝑛−1 ),
then 𝑥𝑛 → 𝑥 ∗ .
BANACH–MAZUR THEOREM: Most well-behaved normed spaces are subspaces of the
space of continuous paths. It is named after Stefan Banach and Stanisław Mazur.
BANACH SPACE: A complete normed vector space on the real or complex numbers.
A. For 1 ≤ 𝑝 < ∞, the space 𝑙𝑝 is a Banach space.
B. The spaces 𝑐0 and 𝑙∞ are Banach spaces.
C. Let 𝐸 be a normed space, and let 𝐹 be a Banach space. Then 𝐵(𝐸, 𝐹) is a Banach
space.
D. Let 1 < 𝑝 < ∞, and again let 𝑞 be such that 1/𝑝 + 1/𝑞 = 1. Then the map
𝑙𝑞 → (𝑙_𝑝)∗ : 𝑢 ↦ 𝜑𝑢 , is an isometric isomorphism, where 𝜑𝑢 is defined, for
∞
𝑢 = (𝑢𝑗 ) ∇ 𝑙𝑞 , by 𝜑𝑢 𝑥 = 𝑖=1 𝑢𝑖 𝑥𝑖 𝑥 = 𝑥𝑖 ∇ 𝑙𝑝

BANACH STAR ALGEBRA: An involution in Banach algebra 𝑅 is an operation 𝑥 → 𝑥 ∗ that

satisfïes
 (𝑥 + 𝑦)∗ = 𝑥 ∗ + 𝑦 ∗ ;
 (𝑙𝑥)∗ = 𝑙𝑥 ∗ ;
 (𝑥𝑦)∗ = 𝑦 ∗ 𝑥 ∗ ;
 (𝑥 ∗ )∗ = 𝑥.
A Banach algebra with an involution is called a Banach *-algebra. A *-homomorphism 𝜑
between two Banach*-algebras is an algebraic homomorphism which preserves
involutions, i.e., 𝜑(𝑥 ∗ ) = 𝜑(𝑥)∗ . To represent a Banach *-algebra, we prefer a *-
representation, i.e., a representation 𝑥 → 𝑇𝑥 on a Hilbert space such that 𝑇𝑥 is equal to
the adjoint 𝑇𝑥 ∗ of 𝑇𝑥 for any 𝑥 ∇ 𝑅 .
BANACH, STEFAN (1892–1945): Stefan Banach was a Polish mathematician who gave a
major contribution to the subject known as functional analysis.
BANACH SUBALGEBRA: A closed subalgebra M of Banach algebra of A is called Banach
subalgebra of A if M itself is a Banach algebra with respect to the operations, identify
and norm defined on A.
BAR GRAPH: A bar graph is a pictorial rendition of statistical data in which
the independent variable can attain only certain discrete values. The dependent
variable may be discrete or continuous. The most common form of bar graph is the
vertical bar graph, also called a column graph. In a vertical bar graph, values of the
independent variable are plotted along a horizontal axis from left to right. Function
values are shown as shaded or colored vertical bars of equal thickness extending
upward from the horizontal axis to various heights. In a horizontal bar graph, the
independent variable is plotted along a vertical axis from the bottom up. Values of the
function are shown as shaded or colored horizontal bars of equal thickness extending
toward the right, with their left ends vertically aligned.

In a specialized type of vertical bar graph called a Pareto chart , values of the dependent
variable are plotted in decreasing order of relative frequency from left to right. Another
type of bar graph called a histogram uses rectangles to show the frequency of data items
in successive numerical intervals of equal size. Other types of bar graphs allow plotting
multiple ranges of the dependent variable, multiple independent variables,
positive/negative variables and multi-category variables.

BARYCENTRIC: Of or related to the centre of mass of an object.

BASES OF VECTOR SPACES: The vectors 𝑣1 , 𝑣2 , ⋯ ⋯ , 𝑣𝑛 in 𝑉 form a basis of 𝑉 if they are

linearly independent and span 𝑉. The vectors 𝑣1 , 𝑣2 , ⋯ ⋯ , 𝑣𝑛 in 𝑉 form a basis of 𝑉 if and
only if ever 𝑣 𝜖 𝑉 can be written uniquely as 𝒗 = 𝛼1 𝑣1 + 𝛼2 𝑣2 + ⋯ ⋯ + 𝛼𝑛 𝑣𝑛 : i.e the
coefficients 𝛼1 , 𝛼2 , ⋯ ⋯ , 𝛼𝑛 are uniquely determined by the vector v.

The scalars 𝛼1 , 𝛼2 , ⋯ ⋯ , 𝛼𝑛 in the note above are called the coordinates of 𝑣 with respect
to the basis 𝑣1 , 𝑣2 , ⋯ ⋯ , 𝑣𝑛 . with respect to the different basis, 𝑣 will have different
coordinates.

Examples: Here are some examples of basis of vector space

1: (1,0) and (0,1) form a basis of ℝ2 .This follows because each element (𝑥, 𝑦) 𝜖ℝ2 can be
written as 𝑥(1,0) + 𝑦(0,1) and this expression is clearly unique.
2: More generally (1,0,0), (0,1,0), (0,0,1) form a basis of ℝ3 .
(1,0,0,0, ), (0,1,0,0), (0,0,1,0) , (0,0,0,1) form a basis of ℝ4 and so on. This is called the
standard basis of ℝ𝑛 for n 𝜖𝑁.

3: There are many other basis of ℝ𝑛 .For example (1,0), (1,1) form a basis of ℝ2 because
(𝛼1 , , 𝛼2 ) = 𝛼1 −𝛼2 1,0 + 𝛼2 (0,1) and this expression is unique. In fact, any two non-
zero vectors such that one is not a scalar multiple of the other form a basis of ℝ2 .

4: As we have defined a basis, it has to consist of a finite number of vectors. Not every
vector space has a finite basis. For example, let 𝐾 𝑥 be the space of polynomials in an
indeterminate x with coefficients in the field K.Let 𝑝1 𝑥 , 𝑝2 𝑥 , … … . , 𝑝𝑛 𝑥 be any finite
collection of polynomials in 𝐾 𝑥 .Then if 𝑑 is the maximum degree of
𝑝1 𝑥 , 𝑝2 𝑥 , … … . , 𝑝𝑛 𝑥 ,any linear combination of 𝑝1 𝑥 , 𝑝2 𝑥 , … … . , 𝑝𝑛 𝑥 has degree
atmost d,and so 𝑝1 𝑥 , 𝑝2 𝑥 , … … . , 𝑝𝑛 𝑥 cannot span 𝐾[𝑥].On the other hand, it is
possible to define what it means for an infinite set of vectors 1, 𝑥, 𝑥 2 , … … … , … . . 𝑥 𝑛 … … ..
is a basis of 𝐾 𝑥 . A vector space with a finite basis is called finite-dimensional. The
spaces of functions mentioned in example above have uncountably infinite dimension.

BASE (TOPOLOGY): Let X be a topological space with Γ as topology defined on it. Then a
collection 𝐵𝑛 of subsets of X is said to be a base for the topology Γ if for every open set
𝐺 ∇ Γ and for every element 𝑥 ∇ 𝐺, there exists a member 𝐵𝑛 of this collection such that
𝑥 ∇ 𝐵𝑛 ⊆ 𝐺. If there exists a countable base for a topological space, then the topological
space is called a second countable space.
BASIC FEASIBLE SOLUTION: A basic feasible solution for a linear programming problem
is a solution that satisfies the constraints of the problem where the number of nonzero
variables equals the number of constraints.
Consider a linear programming problem with 𝑚 constraints and 𝑛 total variables
(including slack variables). Then a basic feasible solution is a solution that satisfies the
constraints of the problem and has exactly 𝑚 nonzero variables and 𝑛 − 𝑚 variables
equal to zero. The basic feasible solutions will be at the corners of the feasible region,
and an important theorem of linear programming states that, if there is an optimal
solution, it will be a basic feasible solution.
BASIC FEASIBLE SOLUTIONS OF TRANSPORTATION PROBLEMS: A feasible solution
(𝑥𝑖,𝑗 ) of a Transportation Problem is said to be basic if there exists a basis 𝐵 for that
Transportation Problem such that 𝑥𝑖,𝑗 = 0 whenever (𝑖, 𝑗) ∈ 𝐵.
BASIC MODULE DECOMPOSITION THEOREM: If 𝑀 is a finitely generated module over a
Euclidean domain 𝐷 then 𝑀 may be written as an internal direct sum 𝑀 = 𝑀1 ⊕
𝑀2 ⊕ · · · ⊕ 𝑀𝑠 where 𝑀𝑖 is a non-trivial cyclic submodule of order 𝑑𝑖 [1 ≤ 𝑖 ≤ 𝑠]
and 𝑑𝑖 |𝑑𝑖+1 [1 ≤ 𝑖 ≤ 𝑠 − 1].
BASIS: A set of vectors in a vector space form a basis if other vectors can be written as a
linear combination of the basis vectors. The vectors in the basis need to be necessarily
linearly independent.
BASIS AND DIMENSION OF A SUBSPACE: A set of vectors 𝑎1 , 𝑎2 , 𝑎3 , … , 𝑎𝑘 belonging to
the subspace 𝑆 is said to be a basic of 𝑆, if

(i) The subspace 𝑆 is spanned by the set 𝑎1 , 𝑎2 , … , 𝑎𝑘 and

(ii) The vectors 𝑎1 , 𝑎2 , … , 𝑎𝑘 are linearly independent.

It can be easily shown that every subspace, 𝑆,or 𝑉𝑛 possesses a basic. It can be easily
shown that the vectors e1 = 1,0,0, … 0 , e2 = 0,1,0, … 0 , e3 = 0,0,1, … 0 , … en =
0,0,0, … 1 constitute a basis of 𝑉𝑛 .

We have already shown that these vectors are linearly independent. Moreover any
vector 𝑎 = (𝑎1 , 𝑎2 , … , 𝑎𝑛 ) of 𝑉𝑛 is expressible as 𝑎 = 𝑎1 𝑒1 + 𝑎2 𝑒2 + 𝑎3 𝑒3 + ⋯ + 𝑎𝑛 𝑒𝑛 .
Hence the vectors 𝑒1 , 𝑒2 , 𝑒3 , … , 𝑒𝑛 constitute a basis of 𝑉𝑛 .

 A basis of a subspace, 𝑆, can always be selected out of a set of vectors which span
𝑆.
 A vector subspace may (and in fact does) possess several bases.
 The number of members in any one basis of a subspace is the same as in any
other basis. This number is called the dimension of the subspace.

BASIS THEOREM: Any linearly independent set of n vectors is a basis of a vector space of
finite dimension n.
BAYES FACTOR: The Bayesian equivalent (in its application) to the classical frequentist
concept of hypothesis testing.
BAYESIAN CONFIDENCE INTERVAL: An interval of a posterior distribution which is
such that the density at any point inside the interval is greater than the density at any
point outside and that the area under the curve for that interval is equal to a
prespecified probability level. For any probability level there is generally only one such
interval, which is also often known as the highest posterior density region. Unlike the
usual confidence interval associated with frequentist inference, here the intervals
specify the range within which parameters lie with a certain probability.

BAYESIAN INFERENCE: An approach to inference based largely on Bayes’ Theorem and

consisting of the following principal steps:

(1) Obtain the likelihood, 𝑓 (𝑥 𝜃 ) describing the process giving rise to the data 𝑥 in
terms of the unknown parameters 𝜃.

(2) Obtain the prior distribution, 𝑓 (𝜃) expressing what is known about 𝜃, prior to
observing the data.

3 Apply Bayes’ theorem to derive the posterior distribution 𝑓 (𝑥 𝜃) expressing what

is known about 𝜃 after observing the data.

(4) Derive appropriate inference statements from the posterior distribution. These may
include specific inferences such as point estimates, interval estimates or probabilities of
hypotheses. If interest centres on particular components of q their posterior
distribution is formed by integrating out the other parameters.

This form of inference differs from the classical form of frequentist inference in several
respects, particularly the use of the prior distribution which is absent from classical
inference. It represents the investigator’s knowledge about the parameters before
seeing the data. Classical statistics uses only the likelihood. Consequently to a Bayesian
every problem is unique and is characterized by the investigator’s beliefs about the
parameters expressed in the prior distribution for the specific investigation.

BAYESIAN PROBABILITY: A popular interpretation of probability which evaluates the

probability of a hypothesis by specifying some prior probability, and then updating in
the light of new relevant data.
BAYES’S RULE PROBABILITY): Bayes’s rule tells how to find the conditional probability
𝑃(𝐵 ∕ 𝐴) (that is, the probability of the event 𝐵 given that event 𝐴 has occurred),
provided that 𝑃(𝐴 ∕ 𝐵) and 𝑃(𝐴 ∕ 𝐵 𝑐 ) are known.
Bayes’s rule states:
𝑃 𝐴 𝐵 𝑃(𝐵)
𝑃 𝐵 𝐴 =
𝑃 𝐴 𝐵 𝑃 𝐵 + 𝑃 𝐴 𝐵 𝑐 𝑃(𝐵 𝑐 )
For example, suppose that two dice are rolled. Let A be the event of rolling doubles, and
let B be the event where the sum of the numbers on the two dice is greater than or equal
to 11. Then 𝑃 𝐴 𝐵 refers to the probability of obtaining doubles if the sum of the two
numbers is greater than or equal to 11; this probability is 1/3. There are 3 possible
outcomes where the sum of the two numbers is greater than or equal to 11, and one of
1
these are doubles: (6, 6). Also, 𝑃 𝐴 𝐵 𝑐 = 1 − 3 = 2/3. Then, we can use Bayes’s rule to

find the probability that the sum of the two numbers will be greater than or equal to 11,
given that doubles were obtained as
1 3
3 12 1
𝑃 𝐵 𝐴 = =
1 3 2 9 7
+
3 12 3 12
BAYES, THOMAS (1702 to 1761): Thomas Bayes was an English mathematician who
studied probability and statistical inference.
BEARING: The angle that the direction makes with north, measured in degrees in a
clockwise direction from north. Bearings of less than 100° are generally written with an
initial 0. For example, north-east has a bearing of 045°, and south-west has a bearing of
225°. Bearing is an important concept in radar, sonar, navigation, and surveying.
BECKMAN–QUARLES THEOREM: If a transformation of the Euclidean plane or a higher-
dimensional Euclidean space preserves unit distances, then it preserves all distances.
Equivalently, every automorphism of the unit distance graph of the plane must be
an isometry of the plane.
BECQUEREL: An SI derived unit of radio acitivity denoted by the symbol Bq.

BEES ALGORITHM: The bees algorithm is a method of problem solving that mimics the
behavior of honeybees to find the optimum solution. Based on the behaviors bees
employ to search and prioritize, the algorithm is a classical example of swarm
intelligence, in which many individuals work together to solve problems or optimize
scenarios. Bees search for food by using scouts to explore areas deemed most likely to
produce favorable results. At first, the scouts conduct random searches to locate the
areas where food exists in the greatest abundance. Then they conduct more orderly,
localized searches until they arrive at the most efficient possible food-recovery process.

The bees algorithm makes it possible for research scientists and engineers to solve
complex problems involving vast amounts of data, categorizing the results according to
specific criteria (such as geographic region or age group), and then giving priority to the
results most likely to yield workable solutions. Computers and swarms of insect
robots can use the bees algorithm as well.

Practical applications of the bees algorithm include:

 Machine vision
 Pattern recognition
 Image analysis
 Job scheduling
 Finding multiple solutions to problem
 Data aggregation
 Mechanical component design
 Robot control
BELL CURVE: It is the shape of the normal curve i.e. the shape of the graph that indicates
a normal distribution in probability and statistics.
BELL-SHAPED DISTRIBUTION: A probability distribution having the overall shape of a
vertical cross-section of a bell. The normal distribution is the most well known example,
but Student’s t-distribution is also this shape.

BELONGS TO: If x is an element of a set S, then we say that x belongs to S and this is
written as x ∇ S. Naturally, x ∈ S means that x does not belong to S.
BELOW (LESS THAN): The limit of a function at 𝑎 from below is the limit of 𝑓(𝑥) as
𝑥 → 𝑎 for values of 𝑥 < 𝑎. It is of particular importance when 𝑓(𝑥) has a discontinuity
at 𝑎, i.e. where the limits from above and from below do not coincide. It can be written
as 𝑓(𝑎−) or lim𝑥→𝑎− 𝑓(𝑥).
BELTRAMI AND ENNPER THEOREM (DIFFERENTIAL GEOMETRY): Statement: At a point
on a surface where the Gaussian curvature is negative and equal to 𝜅, the torsion of the
asymptotic line is ± −𝜅.

BELTRAMI FLOW: If the motion be steady and irrotational then

∂q
∂∅ ∂t = 0 ⇒ curl q = 0. Since 𝑞×curl 𝑞= 0.

It follows that 𝑞 and curl 𝑞 are parallel i.e., streamlines and vortex lines coincide. For
such a motion, 𝑞 is known as a Beltrami vector and the flow is known as a Beltrami flow.

BELTRAMI –MICHELL COMPATIBILITY EQUATIONS:- Let us consider the stress strain

relations

𝜍 1+ 𝜍
ℯ𝑖𝑗 = − 𝛿𝑖𝑗 𝜃 + 𝜏𝑖𝑗
𝐸 𝐸

And the compatibility equation is

ℯ𝑖𝑗 , 𝑘𝑙 + ℯ𝑘𝑙, 𝑖𝑗 − ℯ𝑖𝑘 , 𝑗𝑙 − ℯ𝑗𝑙 , 𝑖𝑘 =0

1 𝜍 1+ 𝜍
Then ∆2 𝜏𝑖𝑗 + 𝜃, 𝑖𝑗 − 1+ 𝜍 𝛿𝑖𝑗 − 𝑑𝑖𝑣 𝐹 = 𝐹𝑖,𝑗 + 𝐹𝑗 ,𝑖
1+ 𝜍 1− 𝜍

1 𝜍
Or ∆2 𝜏𝑖𝑗 + 𝜃, 𝑖𝑗 + 1− 𝜍 𝛿𝑖𝑗 𝑑𝑖𝑣 𝐹 = − 𝐹𝑖,𝑗 + 𝐹𝑗 ,𝑖
1+ 𝜍

1 𝜍
Or ∆2 𝜏𝑖𝑗 + 𝜃, 𝑖𝑗 = − 1− 𝜍 𝛿𝑖𝑗 𝑑𝑖𝑣 𝐹 = − 𝐹𝑖,𝑗 + 𝐹𝑗 ,𝑖
1+ 𝜍

These are called Beltrami- Michell compatibility equations. The Cartesian form of this
equation is

1 𝜕2𝜃 𝜍 2𝜕𝐹𝑥
∆2 𝜏𝑥𝑥 + 2
=− 𝑑𝑖𝑣 𝐹 − ,
1 + 𝜍 𝜕𝑥 1− 𝜍 𝜕𝑥

2
1 𝜕2𝜃 𝜍 2𝜕𝐹𝑦
∆ 𝜏𝑦𝑦 + = − 𝑑𝑖𝑣 𝐹 − ,
1 + 𝜍 𝜕𝑦 2 1− 𝜍 𝜕𝑦
1 𝜕2𝜃 𝜍 2𝜕𝐹𝑧
∆2 𝜏𝑧𝑧 + 2
=− 𝑑𝑖𝑣 𝐹 − ,
1 + 𝜍 𝜕𝑧 1− 𝜍 𝜕𝑧

2
1 𝜕2𝜃 𝜕𝐹𝑦 𝜕𝐹𝑧
∆ 𝜏𝑦𝑧 + =− + ,
1 + 𝜍 𝜕𝑦 𝜕𝑧 𝜕𝑧 𝜕𝑦
 2
1 𝜕2𝜃 𝜕𝐹𝑧 𝜕𝐹𝑥
∆ 𝜏𝑧𝑥 + =− + ,
1 + 𝜍 𝜕𝑧 𝜕𝑥 𝜕𝑥 𝜕𝑧

1 𝜕2𝜃 𝜕𝐹𝑥 𝜕𝐹𝑦

∆2 𝜏𝑥𝑦 + =− +
1 + 𝜍 𝜕𝑥 𝜕𝑦 𝜕𝑦 𝜕𝑥

BELTRAMI’S RESULT: 2n + 1 x 2 − 1 P′n = n n + 1 Pn+1 − Pn−1 .

BELTRAMI'S THEOREM (DIFFERENTIAL GEOMETRY): Geodesic maps preserve the

property of having constant curvature. More precisely, if (𝑀, 𝑔) and (𝑁, 𝑕) are
two Riemannian manifolds and 𝜑 : 𝑀 → 𝑁 is a geodesic map between them, and if
either of the manifolds (𝑀, 𝑔) or (𝑁, 𝑕) has constant curvature, then so does the other
one.
BEND POINT: Also known as an (extremal) turning point.

BERGER–KAZDAN COMPARISON THEOREM: It is a theorem in Riemannian

geometry that gives a lower bound on the volume of a Riemannian manifold and also
gives a necessary and sufficient condition for the manifold to be isometric to the m-
dimensional sphere with its usual metric. The theorem is named after
the mathematicians Marcel Berger and Jerry Kazdan. It can be stated as follows:

Let (𝑀, 𝑔) be a compact 𝑚-dimensional Riemannian manifold with injectivity

radius 𝑖𝑛𝑗(𝑀). Let vol denote the volume form on 𝑀 and let 𝑐𝑚 (𝑟) denote the volume of
the standard m-dimensional sphere of radius r. Then

with equality if and only if (𝑀, 𝑔) is isometric to the 𝑚-sphere 𝑺𝒎 with its usual metric.

BERNOULLI, DANIEL (1700 to 1782): Daniel Bernoulli, son of Johann Bernoulli,

investigated mathematics and other areas. He developed Bernoulli’s theorem in fluid
mechanics, which governs the design of airplane wings.
BERNOULLI DIFFERENTIAL EQUATION: A differential equation which can be
written in the form 𝑦′ + 𝑓(𝑥)𝑦 = 𝑔(𝑥)𝑦 𝑛 . The transformation 𝑧 = 𝑦1−𝑛 gives
𝑑𝑧
= (1 − 𝑛)𝑦 −𝑛 𝑦 ′ . Dividing the original equation by 𝑦 𝑛 gives 𝑦– 𝑛𝑦ʹ + 𝑓(𝑥)𝑦1−𝑛 =
𝑑𝑥
 1 𝑑𝑧
𝑔(𝑥) which reduces to + 𝑓 𝑥 𝑧 = 𝑔(𝑥), which is in linear form and can be
1−𝑛 𝑑𝑥

solved by use of an integrating factor.

BERNOULLI DISTRIBUTION: It is a type of discrete probability distribution. BERNOULLI
EXPERIMENT: An experiment consisting of Bernoulli trials is called a Bernoulli
experiment.

BERNOULLI INEQUALITY: It is an inequality that approximates exponentiations

of 1 + x.

The inequality states that

for every integer 𝑟 ≥ 0 and every real number 𝑥 ≥ − 1. If the exponent 𝑟 is even, then
the inequality is valid for all real numbers 𝑥. The strict version of the inequality reads

for every integer 𝑟 ≥ 2 and every real number 𝑥 ≥ − 1 with 𝑥 ≠ 0.

BERNOULLI, JAKOB (1655 to 1705): Jakob Bernoulli was a Swiss mathematician who
studied concepts in what is now the calculus of variations, particularly the catenary
curve.
BERNOULLI, JOHANN (1667 to 1748): Johann Bernoulli, brother of Jakob Bernoulli, was
also a mathematician investigating these issues.
BERNOULLI NUMBER: Bernoulli number can be defined through the power series
expansion of
𝑥 𝑥 2 ∕ 2! 𝑥 4 ∕ 4! 𝑥 6 ∕ 6!
𝑥 1 − 𝑒 −𝑥 = 1 + + + + +−−−
2 6 30 42
𝑥𝑛
The Bernoulli numbers are the coefficients of in this expansion.
𝑛!

BERNOULLI TRIAL: One of a sequence of independent experiments, each of which has an

outcome considered to be success or failure, all with the same probability p of success.
It is the event that has only two possible outcomes: success and failure.
BERNSTEIN POLYNOMIALS: For 𝑓 ∇ 𝐶[0, 1], the Bernstein polynomials of 𝑓 are given
𝑛 𝑛
by the formula 𝐵𝑛 (𝑓, 𝑥) ∶= 𝑘=0 𝑘 𝑥 𝑘 (1 − 𝑥)𝑛−𝑘 𝑓( 𝑘/𝑛 ). This formula produces a
positive linear map 𝑓 → 𝐵𝑛 (𝑓) of 𝐶[0, 1] into 𝑃𝑛 .
BERTRAND CURVES: If a pair of curves 𝐶 and 𝐶1 are such that the principal normal to 𝐶
are also principal normal to 𝐶1 , then the curves 𝐶 and 𝐶1 are called conjugate or
associate Bertrand curves.

BERTRAND–DIQUET–PUISEUX THEOREM: The theorem is closely related to the Gauss–

Bonnet theorem. Let 𝑝 be a point on a smooth surface 𝑀. The geodesic circle of
radius 𝑟 centered at 𝑝 is the set of all points whose geodesic distance from 𝑝 is equal
to 𝑟. Let 𝐶(𝑟) denote the circumference of this circle, and 𝐴(𝑟) denote the area of the
disc contained within the circle. Then

BERTRAND’S POSTULATE: For any integer greater than 3, there is always at least one
prime between n and 2n – 2. For example, for n = 6 the number 7 is prime and lies
between 6 and 10. This statement was first conjectured in 1845 by Joseph Bertrand.
Bertrand himself verified his statement for all numbers in the interval [2, 3 × 106 ]. His
conjecture was completely proved by Chebyshev in 1852 and so the postulate is also
called the Bertrand–Chebyshev theorem or Chebyshev's theorem. Chebyshev's theorem
can also be stated as a relationship with 𝜋(𝑥), where 𝜋(𝑥) is the prime counting
𝜋
function (number of primes less than or equal to 𝑥): 𝜋 𝑥 − 𝜋 ≥ 1 ∀ 𝑥 ≥ 2.
2

BESSEL, FRIEDRICH WILHELM (1784–1846): Friedrich Wilhelm Bessel was a German

astronomer and mathematician who made a major contribution to mathematics in the
development of what are now called Bessel functions.
BESSEL FUNCTION: These are the functions that provide solutions to Bessel’s
equation. These are of the form
∞
(−1)𝑟 𝑧 𝑛+2𝑟
𝐽𝑛 𝑧 =
𝑟! Γ(𝑛 + 𝑟 − 1) 2
𝑟=0

of which the simplest is the Bessel function of the first kind

∞
(−1)𝑟 𝑧 2𝑟
𝐽0 𝑧 =
𝑟! 2 2
𝑟=0

BESSEL’S EQUATION: It is the second-order differential equation in the form

𝑥 2 𝑦ʺ + (𝑥 2 – 𝑛2 )𝑦 = 0
∂2 V 1 ∂V 1 ∂2 V 𝜕2𝑉
BESSEL EQUATION FROM LAPLACE EQUATION: ∂r 2 + r + r 2 ∂θ 2 + 𝜕𝑧 2 = 0.
∂r
Let V = Rθ′z′ where R, θ′ , z′ are functions of r, θ and z alone, respectively.

Substituting in (i), we have

d2R 1 dR 1 d2θ′ d 2 z′
θ′ z′ dr 2 + r θ′ z′ dr + r 2 Rz′ + Rθ′ dz 2 = 0.
dθ 2

1 d2R 1 dR 1 1 d 2 θ′ 1 d2z′
or +r + r 2 θ′ + z′ + = 0. …………….. ii
R dr 2 dr dθ 2 dz 2

Since the first three terms are independent of z therefore the fourth terms also be
independent of z. Thus the fourth term must be equal to constant

1 d2z′ d2z′
i.e. z′ = C (constant) or = Cz′. ………….. iii
dz 2 dz 2

Similarly third term must be free form θ and therefore must be equal to a constant

1 d2θ′ d2θ′
i.e., = D or = Dθ′ . ………….. iv
θ′ dθ 2 dθ 2

∴ With the help of (iii) and (iv), equation (ii), becomes

d2R dR
r2 + dr + D + Cr 2 R = 0 …………….. v
dr 2

dR dR dv dR d2R d 2 R dv d2R
Putting kr = v so that dr = . dr = k dv and = k dv 2 . dr = k 2 dv 2
dv dr 2

in (v), we have

d2R dR v2
k 2 r 2 dv 2 + kr dv + D + C k 2 R = 0.

d2R dR
Putting 𝐶 = k 2 𝐷 = −n2 , we get v 2 dv 2 + v dv + v 2 − n2 R = 0.

d2y dy
Putting 𝑅 = 𝑦 and 𝑣 = 𝑥, we have x 2 dx 2 + x dx + x 2 − n2 y = 0

which is the Bessel’s equation.

The solution of this equation is called the cylindrical function or Bessel’s function of
order n.
 2 𝑛
BESSEL’S INEQUALITY: If ⌌𝑒𝑖 ⌍ is orthonormal set, then 𝑥 ≥ 𝑖=1 ⌌𝑥, 𝑒𝑖 ⌍ 2 . For an
arbitrary vector 𝑕 of 𝑆, we have

𝑛
2 2
𝑕 ≥ 𝑕, 𝑒𝑖
𝑖=1

BEST APPROXIMATION: A point of a subject in a metric space that is closest to a given

point outside the subset. A best approximation to 𝑓 ∇ 𝑋 from 𝑈 is an element
𝑢∗ ∇ 𝑈 𝑠. 𝑡. 𝑑(𝑓, 𝑢∗ ) = 𝑖𝑛𝑓{𝑑(𝑓, 𝑢) ∶ 𝑢 ∇ 𝑈} = : 𝑑𝑖𝑠𝑡(𝑓, 𝑈).
BETA DISTRIBUTION: If X is a random variable with probability density function given
by
1
𝑓 𝑥, 𝛼, 𝛽 = 𝑥 𝛼 −1 (1 − 𝑥)𝛽−1 ; 0 ≤ 𝑥 ≤ 1
𝐵 𝛼, 𝛽
where 𝐵(𝛼, 𝛽) is a beta function and 𝛼, 𝛽, 𝑥 > 0, then we say that 𝑋 has a beta
𝛼
distribution with parameters 𝛼, 𝛽. The beta distribution has mean and variance
𝛼+𝛽
𝛼𝛽
2
.
𝛼 +𝛽 𝛼+𝛽+1

BETA FUNCTION: The function

1
Γ 𝑝 Γ(𝑞)
𝐵 𝑝, 𝑞 = 𝑥 𝑝−1 1 − 𝑥 𝑞−1
𝑑𝑥 =
Γ(𝑝 + 𝑞)
0

where 𝛤(𝑝) is the gamma function, is called a beta function. For integers 𝑚, 𝑛, we have
𝑚 − 1 ! (𝑛 − 1)!
𝛽 𝑚, 𝑛 =
(𝑚 + 𝑛 − 1)!

BÉZOUT'S IDENTITY: Let 𝑎 and 𝑏 be nonzero integers and let 𝑑 be their greatest
common divisor. Then there exist integers 𝑥 and 𝑦 such that

In addition,

 𝑑 is the smallest positive integer that can be written as 𝑎𝑥 + 𝑏𝑦.

 every integer of the form 𝑎𝑥 + 𝑏𝑦 is a multiple of 𝑑.

𝑥 and 𝑦 are called Bézout coefficients for (𝑎, 𝑏) and are not unique. A pair of Bézout
coefficients can be computed by the extended Euclidean algorithm.
BHĀSKARA (1114–85): Bhaskara was an eminent Indian mathematician who continued
in the tradition of Brahmagupta, making corrections and filling in many gaps in the
earlier work. He solved examples of Pell’s equation and grappled with the problem of
division by zero.

BHATIA–DAVIS INEQUALITY: It is named after Rajendra Bhatia and Chandler Davis, is

an upper bound on the variance of any bounded probability distribution on the real line.

Suppose a distribution has minimum 𝑚, maximum 𝑀, and expected value 𝜇. Then the
inequality says:

Equality holds precisely if all of the probability is concentrated at the

endpoints 𝑚 and 𝑀.

Bi: It is the abbreviation for Airy function of the second kind.

BIAS: A prejudice or a lack of objectivity or randomness resulting in an imbalance that
makes it likely that the outcome will tend to be distorted. In statistics this is when a
process contains some systematic imbalance so that, on average, the outcome of the
process is not equal to the true value. Randomization techniques are employed to try
and remove bias that may result from other sample estimator selection methods
requiring choices to be made. An estimate of a parameter is biased if its expected value
does not equal the population value of the parameter.
BIASED SAMPLE: It is a sample whose composition is not determined only by the
population from which it has been taken, but also by some property of the sampling
method which has a tendency to cause an over-representation of some parts of the
population. It is a property of the sampling method rather than the individual sample.
BICONDITIONAL STATEMENT: A biconditional statement is a compound statement that
says one sentence is true if and only if the other sentence is true. Symbolically, this is
written as p ↔ q, which means “p → q” and “q → p.”
BIEBERBACH CONJECTURE: Because of applications, there is a lot of interest in
functions 𝑓 which are analytic and one-one in 𝐵(0, 1). Since 𝑓 ′ (0) ≠ 0, we can
normalize so that 𝑓(0) = 0, 𝑓 ′ (0) = 1 (else consider (𝑓(𝑧) − 𝑓(0))/𝑓 ′ (0) instead of
∝
𝑓), and look at the Taylor series 𝑓(𝑧) = 𝑧 + 𝑘=2 𝑎𝑘 𝑧 𝑘 , |𝑧| < 1. Bieberbach proved
in 1916 that |𝑎2 | ≤ 2, and asked whether we always have |𝑎𝑛 | ≤ 𝑛. This became his
conjecture, an object of an enormous amount of research, and was finally proved by
Louis de Branges in 1984
BIHOLOMORPHISM: 𝑓 ∶ 𝐶 → 𝐶 is said to be a biholomorphism if it is bijective and
𝑓, 𝑓 −1 are both holomorphic. Since the derivative map is a composition of scaling and
rotation, it preserves angles between vectors. So biholomorphisms are conformal maps,
meaning they preserve angles.
BIJECTION: A one-to-one onto mapping, that is, a mapping that is both injective and
surjective is known as a bijection.
BILATERAL SYMMETRY: Reflective symmetry. It is a type of self-similarity
through reflection in a line or a plane.

BILINEAR TRANSFORMATION ON MOBIUS TRANSFORMATION: A transformation of the

form

𝑎𝑧 + 𝑏
𝑤=
𝑐𝑧 + 𝑑

is called a bilinear transformation of linear fractional transformation, where 𝑎, 𝑏, 𝑐, 𝑑 are

complex constants. Such type of transformation was first studied by Mobius and hence
it is sometimes called Mobius transformation is expressible as

𝑐𝑤𝑧 + 𝑤𝑑 − 𝑎𝑧 − 𝑏 = 0

Evidently it is linear both in 𝑤 𝑎𝑛𝑑 𝑧.

That is why, it is called bilinear. Here we assume that 𝑎𝑑 − 𝑏𝑐 ≠ 0 which is called the
determinant of the transformation. The transformation (1) is said to be normalized if
𝑎𝑑 − 𝑏𝑐 = 1. If the determinant vanishes, then 𝑤 is merely a constant.

BI-LIPSCHITZ MAP: A map 𝑓: 𝑋 → 𝑌 is called bi-Lipschitz if there are positive constants c

and C such that for any x and y in X

BILLION: The name variously given to a thousand million and a million million in the
past. Traditionally, the US used a billion as the former and the UK as the latter, also
referred to as short billion and long billion respectively. With the rapid globalization in
the past few decades, the US definition has become more prevalent although it is better
to clarify as its use is dominant.
BIMODAL: A frequency distribution of numerical data that shows two distinct peaks
(modes).
BINARY AND TERNARY CODE: A code over 𝐴 = {0, 1} is called a binary code and a code
over 𝐴 = {0, 1, 2} is called a ternary code.
BINARY HAMMING CODE: Let 𝑟 ≥ 2 and let C be a binary linear code with 𝑛 = 2𝑟 −
1 whose parity-check matrix 𝐻 is such that the columns are all of the non-zero vectors
in 𝐹2𝑟 . This code 𝐶 is called a binary Hamming code of length 2𝑟 − 1, denoted
𝐻𝑎𝑚(𝑟, 2).
BINARY NUMBERS: Binary (base-2) numbers are written in a positional system that
uses only two digits: 0 and 1. Each digit of a binary number represents a power of 2. The
rightmost digit is the 1’s digit, the next digit to the left is the 2’s digit, and so on. For
example, the binary number 11101 represents
1 × 24 + 1 × 23 + 1 × 22 + 0 × 21 + 1 × 20 = 16 + 8 + 4 + 0 + 1 = 29
BINARY OPERATIONS ON SETS: A binary operation ∗ on a set 𝐴 is an operation which,
when applied to any elements 𝑥 and 𝑦 of the set 𝐴, yields an element x ∗ 𝑦 of 𝐴. Example
The arithmetic operations of addition, subtraction and multiplication are binary
operations on the set 𝑅 of real numbers which, when applied to real numbers 𝑥 and 𝑦,
yield the real numbers 𝑥 + 𝑦, 𝑥 − 𝑦 and 𝑥𝑦 respectively. However division is not a
binary operation on the set of real numbers, since the quotient 𝑥/𝑦 is not defined when
𝑦 = 0. (Under a binary operation ∗ on a set must determine an element 𝑥 ∗ 𝑦 of the set
for every pair of elements 𝑥 and 𝑦 of that set.)
BINARY RELATIONS: A binary relation on a set specifies relations between pairs of
elements from the set. Example The relations = ‘equals’ , ≠ ‘not equal to’ , < ‘less
than’ , > ‘greater than’ , ≤ ‘less than or equal to’ and ≥ ‘greater than or equal to’
are all binary relations on the set 𝑅 of real numbers. Let 𝐴 be a set, and let 𝑃(𝐴) be the
power set of 𝐴 (i.e., the set whose elements are the subsets of 𝐴). Then ⊂ is a binary
relation on 𝑃(𝐴), where two subsets 𝐵 and 𝐶 of A satisfy 𝐵 ⊂ 𝐶 if and only if 𝐵 is a
subset of 𝐶. If one has a relation 𝑅 on a set 𝐴, then, given two elements 𝑥 and 𝑦 of 𝐴,
either 𝑥 is related to 𝑦, in which case we may write 𝑥𝑅𝑦, or else the element 𝑥 is not
related to 𝑦.
BINARY SYMMETRIC CHANNEL: A binary symmetric channel has two probabilities:
𝑃[1 𝑟𝑒𝑐𝑒𝑖𝑣𝑒𝑑 | 0 𝑤𝑎𝑠 𝑠𝑒𝑛𝑡] = 𝑃[0 𝑟𝑒𝑐𝑒𝑖𝑣𝑒𝑑 | 1 𝑤𝑎𝑠 𝑠𝑒𝑛𝑡] = 𝑝
𝑃[1 𝑟𝑒𝑐𝑒𝑖𝑣𝑒𝑑 | 1 𝑤𝑎𝑠 𝑠𝑒𝑛𝑡] = 𝑃[0 𝑟𝑒𝑐𝑒𝑖𝑣𝑒𝑑 | 0 𝑤𝑎𝑠 𝑠𝑒𝑛𝑡] = 1 − 𝑝
The probability 𝑝 is called the crossover probability.

BINARY SYSTEM: A base 2 positional system representing numbers with only

the digits 0 and 1.

BINARY VARIABLE: A random variable with only two possible outcomes.

BINOMIAL: A binomial is the sum of two terms. For example, (𝑝𝑥 + 𝑞) is a binomial.

BINOMIAL COEFFICIENTS: Numbers used in the calculation of the coefficients of

binomial expansions. They can be represented in Pascal's triangle and are identical to

the concept of combinations in combinatorics. The symbol is used in the calculation

of the (𝑘 + 1)st term in a binomial expansion of exponent 𝑛, it is also the number of
ways that a set of 𝑘 objects can be chosen from a set of n distinct objects, provided that
the set of 𝑘 objects are chosen at the same time and so the order (of being chosen) does
not matter. The calculation of these binomial coefficients involves factorials as follows,

Note that the following is true for all values of 𝑛 ≥ 0

BINOMIAL DISTRIBUTION: Suppose we want to conduct an experiment 𝑛 times, with a

probability of success of 𝑝 each time. If X is the number of successes that occur in those
𝑛 trials, then X will have the binomial distribution with parameters 𝑛 and 𝑝. X is a
discrete random variable whose probability function is given by
𝑛 𝑟
𝑃 𝑋=𝑟 = 𝑝 (1 − 𝑝)𝑛−𝑟
𝑟
 The expectation is 𝐸 𝑋 = 𝑛𝑝
 The variance is 𝑉𝑎𝑟 𝑋 = 𝑛𝑝(1 − 𝑝).

BINOMIAL INVERSE THEOREM: If 𝑨, 𝑼, 𝑩, 𝑽 are matrices of sizes 𝑝 × 𝑝, 𝑝 × 𝑞, 𝑞 ×

𝑞, 𝑞 × 𝑝 respectively, then
provided A and B + BVA−1UB are nonsingular. Note that if B is invertible, the
two B terms flanking the quantity inverse in the right-hand side can be replaced with
(B−1)−1, which results in

BINOMIAL THEOREM: The binomial theorem tells how to expand the expression
(𝑎 + 𝑏)𝑛 . The general formula is
𝑛 0 𝑛 1 𝑛 𝑛
(𝑎 + 𝑏)𝑛 = 𝑝 (1 − 𝑝)𝑛 + 𝑝 (1 − 𝑝)𝑛−1 + − − − + 𝑝 (1 − 𝑝)0
0 1 𝑛
Some examples of the powers of binomials are as follows:
(𝑎 + 𝑏)0 = 1
(𝑎 + 𝑏)1 = 𝑎 + 𝑏
(𝑎 + 𝑏)2 = 𝑎2 + 2𝑎𝑏 + 𝑏 2
(𝑎 + 𝑏)3 = 𝑎3 + 3𝑎2 𝑏 + 3𝑎𝑏 2 + 𝑏 3
(𝑎 + 𝑏)4 = 𝑎4 + 4𝑎3 𝑏 + 6𝑎2 𝑏 2 + 4𝑎𝑏 3 + 𝑏 4
(𝑎 + 𝑏)5 = 𝑎5 + 5𝑎4 𝑏 + 10𝑎3 𝑏 2 + 10𝑎2 𝑏 3 + 5𝑎𝑏 4 + 𝑏 5
The coefficients form an interesting pattern of numbers known as Pascal’s triangle. This
triangle is an array of numbers such that any entry is equal to the sum of the two entries
above it.
BINORMAL: The normal perpendicular to the principal to the principal normal at point
𝑃 is called binormal at point 𝑃.
BIPARTITE GRAPHS: A graph (𝑉, 𝐸) is said to be bipartite if there exist subsets 𝑉1 and 𝑉2 ,
such that

(i) 𝑉1 ∪ 𝑉2 = 𝑉 ;
(ii) 𝑉1 ∩ 𝑉2 = ∅;
(iii) each edge in E is of the form {𝑣, 𝑤} with 𝑣 ∇ 𝑉1 and 𝑤 ∇ 𝑉2 .

BIQUADRATIC: A polynomial that is the quadratic of a square. The result is

a quartic whose cubic and linear coefficients are zero. Note that this term is also used to
refer the product of two quadratics or the quadratic of a quadratic, which makes
the term equivalent to the term quartic.

BIRCH AND SWINNERTON-DYER CONJECTURE: The Birch and Swinnerton-Dyer

conjecture on elliptic curves postulates a connection between the rank of an elliptic
curve and the order of pole of its Hasse-Weil L-function. It has been an significant
pointer in Diophantine geometry since the mid-1960s, with significant results such as
the Coates–Wiles theorem, Gross–Zagier theorem and Kolyvagin's theorem.
BIRTHDAY PROBLEMS: The birthdays of 𝑟 people form a sample of size 𝑟 from the
population of all days in the year. As a first approximation these may be considered as a
random selection of birthdays from the year consisting of 365 days. Then the
probability, 𝑝, that all 𝑟 birthdays are different is

1 2 3 r−1
𝒫 = 1− 1− 1− −−− 1−
365 365 365 365

For example, when 𝑟 = 23, 𝑝 < 0.5, so that for 23 people the probability that no two
people have the same birthday is less than a half, so the probability that at least two of
the twenty three people share a birthday is greater than a half; most people when asked
to make a guess of how many people are needed to achieve a greater than 50% chance
of at least two of them sharing a birthday, put the figure higher than 23. Matching
triplets of birthdays are much harder to observe and it can be shown that in a group of
23 people the probability of at least one triplet of matching birthdays is only 0.013; now
there has to be a group of 88 people before the probability becomes larger than 0.5.

BIRECTANGULAR: A plane figure with two right angles.

BISECT: Dividing into two equal halves.

BISECTOR: A line or plane which divides a geometrical object (e.g. a figure or an angle)
into two congruent parts.

BIT: The simplest unit of information consisting of one binary digit, conventionally
labelled 0 and 1.

BITANGENT: The tangent of a curve simultaneously at two different points.

BLASIUS’S THEOREM: Let a fixed cylinder be placed in a liquid which is moving steadily
and irrotationally given by the relation 𝑤 = 𝑓(𝑧), if the hydro-dynamical pressure on
the contour of a fixed cylinder are represented by a force (𝑋, 𝑌) and a couple 𝑀 about
the origin of coordinates then
 1 𝑑𝑤 2
𝑋 − 𝑖𝑌 = 2 𝑖𝜌 ∫𝐶 𝑑𝑧 (neglecting extraneous forces)
𝑑𝑧

1 𝑑𝑤 2
and 𝑀 = 𝑅𝑒𝑎𝑙 𝑝𝑎𝑟𝑡 𝑜𝑓 − 2 𝜌 ∫𝐶 𝑑𝑧 ,
𝑑𝑧

where the integrations are round any contour which surrounds the cylinder.

BLOCH'S THEOREM: Bloch's theorem gives a lower bound on the size of a disc in which
an inverse to a holomorphic function exists. Let 𝑓 be a holomorphic function in the unit
disk |𝑧| ≤ 1. Take |𝑓′(0)| = 1. Then there exists a disc of radius 𝑏 and an analytic
function 𝜑 in this disc, such that 𝑓(𝜑(𝑧)) = 𝑧 for all 𝑧 in this disc. Here 𝑏 > 1/72 is an
absolute constant.

BLOCK CODE: Let 𝐴 = {𝑎1 , . . . , 𝑎𝑞 } be an alphabet; we call the 𝑎𝑖 values symbols. A block
code 𝐶 of length 𝑛 over 𝐴 is a subset of 𝐴𝑛 . A vector 𝑐 ∇ 𝐶 is called a codeword. The
number of elements in 𝐶, denoted |𝐶|, is called the size of the code. A code of length 𝑛
and size 𝑀 is called an (𝑛, 𝑀)-code.
BOCHNER’S THEOREM: A continuous function 𝑓 ∶ 𝐺 → 𝐶 is positive definite if and only
if 𝑓 = µ for some µ ∇ 𝑀 +(𝐺 ).

BODY FORCES: The forces which are proportional to the mass contained in the volume
element are called body forces.

BOLYAI, JANOS (1802 to 1860): Janos Bolyai was a Hungarian mathematician who
developed a version of non-Euclidian geometry.
BOLZANO WEIRSTRASS THEOREM: Every infinite bounded set of real numbers has a
limit point.
Every bounded sequence has at least one limit point.
BONNET’S THEOREM ON PARALLEL SURFACE DIFFERENTIAL GEOMETRY): In general
for every surface (S) of constant positive Gaussian curvature. 𝐴−2 there are associated
two surface of constant mean curvature ±(2𝐴)−1 which are parallel to the former (S)
and distant ± 𝐴 form it.

BONSE'S INEQUALITY: if 𝑝1 , 𝑝2 , . . . , 𝑝𝑛+1 are the smallest 𝑛 + 1 prime

numbers and 𝑛 ≥ 4, then
BOOLE GEORGE (1815 to 1865): George Boole was an English mathematician who
developed the symbolic analysis of logic now known as Boolean algebra, which is used
in the design of digital computers.
BOOLEAN ALGEBRA: Boolean algebra is the study of operations carried out on variables
that can have only two values: 1 (true) or 0 (false). Boolean algebra was developed by
George Boole in the 1850s; it is an important part of the theory of logic and has become
of tremendous importance since the development of computers.
Here are some rules from Boolean algebra. In the following statements, p, q, and r
represent Boolean variables and ⟷ represents “is equivalent to.” Parentheses are used
as they are in arithmetic: an operation inside parentheses is to be done before the
operation outside the parentheses.
Double Negation:
𝑝 ↔ 𝑁𝑂𝑇 (𝑁𝑂𝑇 𝑝)
Commutative Principle:
(𝑝 𝐴𝑁𝐷 𝑞) ↔ (𝑞 𝐴𝑁𝐷 𝑝)
(𝑝 𝑂𝑅 𝑞) ↔ (𝑞 𝑂𝑅 𝑝)
Associative Principle:
𝑝 𝐴𝑁𝐷 (𝑞 𝐴𝑁𝐷 𝑟) ↔ (𝑝 𝐴𝑁𝐷 𝑞) 𝐴𝑁𝐷 𝑟
𝑝 𝑂𝑅 (𝑞 𝑂𝑅 𝑟) ↔ (𝑝 𝑂𝑅 𝑞) 𝑂𝑅 𝑟
Distribution:
𝑝 𝐴𝑁𝐷 (𝑞 𝑂𝑅 𝑟) ↔ (𝑝 𝐴𝑁𝐷 𝑞) 𝑂𝑅 (𝑝 𝐴𝑁𝐷 𝑟)
𝑝 𝑂𝑅 (𝑞 𝐴𝑁𝐷 𝑟) ↔ (𝑝 𝑂𝑅 𝑞) 𝐴𝑁𝐷 (𝑝 𝑂𝑅 𝑟)
De Morgan’s Laws:
(𝑁𝑂𝑇 𝑝) 𝐴𝑁𝐷 (𝑁𝑂𝑇 𝑞) ↔ 𝑁𝑂𝑇 (𝑝 𝑂𝑅 𝑞)
(𝑁𝑂𝑇 𝑝) 𝑂𝑅 (𝑁𝑂𝑇 𝑞) ↔ 𝑁𝑂𝑇 (𝑝 𝐴𝑁𝐷 𝑞)

BOOLEAN FUNCTION: Let 𝐴 be a set. A Boolean function on 𝐴 is a function 𝑓: 𝐴 → {𝑇, 𝐹}

whose domain is 𝐴 and whose codomain is the set {𝑇, 𝐹} whose elements are the truth
values T = true and F = false.

BOOLEAN PRIME IDEAL THEOREM: The Boolean prime ideal theorem is the strong
prime ideal theorem for Boolean algebras. Its formal statement is:
Let 𝐵 be a Boolean algebra, let 𝐼 be an ideal and let 𝐹 be a filter of 𝐵, such
that 𝐼 and 𝐹 are disjoint. Then 𝐼 is contained in some prime ideal of 𝐵 that is disjoint
from 𝐹.

The weak prime ideal theorem for Boolean algebras simply states that every Boolean
algebra contains a prime ideal.

BOOLE'S INEQUALITY: For any finite or countable set of events, the probability that at
least one of the events happens is no greater than the sum of the probabilities of the
individual events. Formally, for a countable set of events 𝐴1 , 𝐴2 , 𝐴3 , . . ., we have

BOOTSTRAP: A statistical re-sampling method for obtaining estimators of

distribution parameters.

BOREL AND CARATHEODORY THEOREM (COMPLEX ANALYSIS): Suppose 𝑓 𝑧 is

analytic for 𝑧 ≤ 𝑅. Also suppose 𝑀 𝑟 𝑎𝑛𝑑 𝐴(𝑟) are respectively maxima of 𝑓(𝑧) and
𝑅 𝑓(𝑧) on 𝑧 =r where 0 < 𝑟 < 𝑅. then

2𝑟 𝑅+𝑟
𝑀 𝑟 ≤ 𝐴 𝑅 + 𝑓(0)
𝑅−𝑟 𝑅−𝑟

BOREL, FÉLIX EDOUARD JUSTIN EMILE: Félix Edouard Justin Emile Borel (1871–1956)
was a French mathematician who was one of the first to study real-valued functions,
with important results in set theory and measure theory.
BOREL MEASURE: It is a measure defined on the sigma algebra of a topological space
onto the set of real numbers. If the mapping is onto the interval [0, 1] it is a Borel
probability measure.
BOREL SET: It is a set obtained from repeated applications of unions and intersections
of countable collections of closed or open intervals on the real line.
BOREL’S METHOD OF SUMMATION: If for a given series 𝑢𝑛 ,

∞
𝑠𝑛 𝑥 𝑛
𝑢 𝑥 =
𝑛!
𝑛=0
is convergent for all 𝑥, and 𝑢(𝑥) 𝑒 𝑥 → 𝑠 𝑎𝑠 𝑥 → ∞, then 𝑢𝑛 is said to be summable by
Borel’s expontial method to the sum 𝑠, and we write 𝑢𝑛 = 𝑠(𝐵). The transformation
thus determined is denoted by 𝐵 and is called Borel’s method of summation.

BOREL 𝝈 −ALGEBRA: Let 𝐾 be a compact topological space. The Borel 𝜍 −algebra, 𝐵(𝐾),
on 𝐾, is the 𝜍 −algebra generated by the open sets in 𝐾. A member of 𝐵(𝐾) is a Borel
set. Notice that if 𝑓: 𝐾 → 𝐾 is a continuous function, then clearly f is 𝐵(𝐾) −measurable
(the inverse image of an open set will be open, and hence certainly Borel). So if
µ: 𝐵(𝐾) → 𝐾 is a finite charge or complex measure (for 𝐾 = ℝ or 𝐾 = ℂ respectively),
then 𝑓 will be µ-integrable (as 𝑓 is bounded) and so we can define 𝜑𝜇 : 𝐶𝐾 𝐾 → 𝐾 by
𝜑𝜇 𝑓 = ∫𝑋 𝑓𝑑𝜇. Clearly 𝜑𝜇 is linear. Suppose for now that µ is positive, so that
𝜑𝜇 (𝑡) ≤ ∫𝑋 𝑓 𝑑𝜇 ≤ 𝑓 ∞ 𝜇(𝐾)

So 𝜑𝜇 ∇ 𝐶𝐾 (𝐾)∗ with ||𝜑𝜇 || ≤ µ(𝐾).

A measure µ: 𝐵(𝐾) → [0, ∞) is regular if for each 𝐴 ∇ 𝐵(𝐾), we have

𝑠𝑢𝑝 𝜇 𝐸 : 𝐸 ⊆ 𝐴 𝑎𝑛𝑑 𝐸 𝑖𝑠 𝑐𝑜𝑚𝑝𝑎𝑐𝑡

𝜇 𝐴 =
𝑖𝑛𝑓 𝜇 𝐹 : 𝐴 ⊆ 𝐹 𝑎𝑛𝑑 𝐹 𝑖𝑠 𝑐𝑜𝑚𝑝𝑎𝑐𝑡

A charge 𝜈 = 𝜈+ − 𝜈− is regular if 𝜈+ and 𝜈− are regular measures. A complex measure is

regular if its real and imaginary parts are regular.

Note that

1. Many common measures on the real line, e.g. the Lebesgue measure, point
measures, etc., are regular.

2. Examples of measure µ on [0,1] which are not regular:

1
𝜇 ∅ = 0, 𝜇 = 0, 𝜇 𝐴 = +∞
2

3. Another example of a 𝜍 −additive measure µ on [0,1] which is not regular:

0, 𝑖𝑓 𝐴 𝑖𝑠 𝑎𝑡 𝑚𝑜𝑠𝑡 𝑐𝑜𝑢𝑛𝑡𝑎𝑏𝑙𝑒
𝜇 𝐴 =
+∞, 𝑖𝑓 𝐴 𝑖𝑠 𝑜𝑡𝑕𝑒𝑟𝑤𝑖𝑠𝑒

We let 𝑀𝑅 (𝐾) and 𝑀𝐶 (𝐾) be the collection of all finite, regular, signed or complex
measures on 𝐵(𝐾). The variation || · || is a norm on 𝑀𝐾 (𝐾).
BORSUK-ULAM THEOREM: For any continuous real function on the sphere 𝑆 𝑛 there
must be antipodal (i.e. opposite) points where the values of the function are same.

BOUND: Let 𝑆 be a non-empty subset of 𝑹. A real number 𝑢 is said to be an upper bound

for 𝑆 if 𝑢 is greater than or equal to every element of 𝑆. If 𝑆 has an upper bound, then 𝑆 is
said to be a set bounded above. Moreover, 𝑢 is a supremum (or least upper bound) of 𝑆
if 𝑢 is an upper bound for 𝑆 and no upper bound for 𝑆 is less than 𝑢; this is written
𝑢 = 𝑠𝑢𝑝 𝑆. For example, if 𝑆 = {1/𝑛; 𝑛 ∇ 𝑁} then 𝑠𝑢𝑝 𝑆 = 1. Similarly, a real number
𝑙 is a lower bound for the set 𝑆 if 𝑙 is less than or equal to every element of 𝑆. If 𝑆 has a
lower bound, then 𝑆 is said to be a set bounded below. Moreover, 𝑙 is an infimum (or
greatest lower bound) of 𝑆 if 𝑙 is a lower bound for the set 𝑆 and no lower bound for 𝑆 is
greater than 𝑙; this is written 𝑙 = 𝑖𝑛𝑓 𝑆. A set is said to be a bounded set if it is bounded
above and below. It is a non-elementary result about the real numbers that any non-
empty set that is bounded above has a supremum, and any non-empty set that is
bounded below has an infimum.
BOUNDARY CONDITION: A complete set of values for all variables at some instant (often
the initial condition, when 𝑡 = 0) which provides a particular solution to a differential
equation.
BOUNDARY MAXIMUM MODULUS THEOREM: Let 𝑈 ⊆ 𝐶 be a bounded domain. Let 𝑓 be
a continuous function on 𝑈 that is holomorphic on 𝑈. Then the maximum value of |𝑓| on
𝑈 must occur on ∂U.
BOUNDARY MINIMUM MODULUS PRINCIPLE: Let 𝑈 ⊆ 𝐶 be a bounded domain. Let 𝑓 be
a continuous function on 𝑈 that is holomorphic on 𝑈. Assume that 𝑓 never vanishes on
𝑈. Then the minimum value of |𝑓| on 𝑈 must occur on ∂U.
BOUNDARY VALUE PROBLEM: A differential equation to be satisfied over a region
together with a set of boundary conditions is usually called a boundary value problem.
BOUNDED FUNCTION: A real valued function 𝑓, defined on a domain 𝐷, is said to be a
bounded function on 𝐷 if there is a number 𝐾 such that, for all 𝑥 in 𝐷, we have
| 𝑓(𝑥)| < 𝐾. If 𝑓 is a continuous function on a closed interval [𝑎, 𝑏] then it is bounded
on 𝑎, 𝑏 .
BOUNDED INVERSE THEOREM: Bounded inverse theorem is an important result in the
theory of bounded linear operators on Banach spaces. It states that a bijective bounded
linear operator 𝑇 from one Banach space to another has bounded inverse 𝑇 −1 . It is
equivalent to both the open mapping theorem and the closed graph theorem.

BOUNDED LINEAR FUNCTIONAL: The linear functional Φ on the normed linear space 𝑁
is said to be bounded if there exists a real constant 𝑘, such that Φ (𝑓) ≤ 𝑘 𝑓 .

Equivalently, we can say that Φ is bounded if Φ (𝑓) is bounded on the closed unit
sphere of 𝑁 i.e, Φ(𝑓) / 𝑓 is bounded for every 𝑓 ≠ 0 in 𝑁.

BOUNDED LINEAR OPERATOR: Let 𝑁 and 𝑁′ be two normed linear spaces with the
same scalars. Then a linear operator 𝑇 of 𝑁 into 𝑁′ is said to be bounded if there exists
a non-negative real number 𝑘 such that

𝑇 𝑓 ≤ 𝑘 𝑓 , for all 𝑓 ∇ 𝑁.

BOUNDED SEQUENCE: The real sequence ⌌𝑎𝑛 ⌍ is said to be bounded if there is a real
number 𝑀 such that, for all 𝑛, |𝑎𝑛 | < 𝑀.
BOUNDED SET: A set 𝐴 is said to be a bounded set if closure of 𝐴 does not consist of any
infinite values.

BOX-AND-WHISKER DIAGRAM: (box plot) A graphical summary of data representing the

boundaries of the 4 quadrants on a scale with the middle quadrants represented
as rectangles (boxes) and the outer quadrants as lines (whiskers).

BOX PLOT: Box plot is another name of the box-and-whisker diagram, also known as a
box-and-whisker plot.

BRAHMAGUPTA (598–665): Brahmagupta was an Indian astronomer and

mathematician whose text on astronomy includes some notable mathematics for its
own sake: the areas of quadrilaterals and the solution of certain Diophantine equations,
for example. Here, the systematic use of negative numbers and zero occurs for probably
the first time.

BRACES: Symbols { and } commonly used to represent the order of operations along
with parentheses and brackets. (And other less common types.) They're also used to
represent sets.
BRACHISTOCHRONE: The trajectory of fastest travel between 2 points if a bead is
considered to be guided by a smooth wire, under gravity only, between the 2 points
starting from rest. Note that the starting point must not be the lower of the 2 points or
the bead would never reach the other point. (At the starting point, the bead
has zero Kinectic energy. It will never reach the point with higher gravitational potential
energy. However, due tot he smoothness of the wire, points of equal heights are
possible.)

BRACKETS: Symbols [ and ] commonly used to represent the order of operations, along
with parentheses and braces. (And other less common types.) Informally, all 3 pairs of
symbols are referred to as brackets. Brackets are also used to
represent matrices amongst other mathematical objects.

BRANCH: A branch is a section of a curve with an endpoint at which it meets another

branch, and where the differential has a discontinuity.
BRANCH OF A HYPERBOLA: The two separate parts of a hyperbola are called the two
branches.
BRANCH AND BOUND METHOD: This procedure constructs a branching method that
terminates each branch once the constraint limit has been reached, and by allowing
items to be added in an order which is determined at the start. This reduces
considerably the number of combinations that have to be tried.
BRANCH POINT: Branch point is a point on the curve at which two or more branches of
the curve meet.
BRANCH POINT (COMPLEX ANALYSIS): A point is called a branch point of a function
𝑓 𝑧 if some of the branches interchange as the independent variable 𝑧 describes a
closed path about it.

BRITISH RAIL EXPRESS METRIC: If we define 𝑑 ∶ 𝑅 2 × 𝑅 2 → 𝑅 by

2
𝑢 + 𝑣 2 , 𝑖𝑓 𝑢 ≠ 𝑣,
𝑑(𝑢, 𝑣) =
0 𝑖𝑓 𝑢 = 𝑣,

then 𝑑 is called the British Rail express metric.

BROOKS' THEOREM: It states a relationship between the maximum degree of a graph

and its chromatic number. According to the theorem, in a connected graph in which
every vertex has at most 𝑚 neighbors, the vertices can be colored with only 𝑚 colors,
except for two cases, complete graphs and cycle graphs of odd length, which require
𝑚 + 1 colors.
BROUWER'S FIXED-POINT THEOREM: It is a fixed-point theorem in topology, named
after Luitzen Brouwer. It states that for any continuous function 𝑓 mapping a compact
convex set into itself there is a point 𝑥0 such that 𝑓(𝑥0 ) = 𝑥0 . The simplest forms of
Brouwer's theorem are for continuous functions 𝑓 from a closed interval 𝐼 in the real
numbers to itself or from a closed disk 𝐷 to itself.

BRUNN–MINKOWSKI THEOREM: It is a theorem relating the Lebesgue measures

of compact subsets of Euclidean space. Let 𝑛 ≥ 1 and let 𝜇 denote the Lebesgue
measure on Rn. Let 𝐴 and 𝐵 be two nonempty compact subsets of Rn. Then the
following inequality holds:

where 𝐴 + 𝐵 denotes the Minkowski sum:

BUMP FUNCTION: Let 𝜍: 𝑈 → 𝑀 be a chart on an abstract manifold, and let

𝐵(𝑥0 , 𝑟) ⊂ 𝐵 (𝑥0 , 𝑠) ⊂ 𝑈 be concentric balls in 𝑈 with 0 < 𝑟 < 𝑠. There exists a
smooth function 𝑔 ∇ 𝐶 ∞ (𝑀), which takes values in [0, 1], such that 𝑔(𝑞) = 1 for
𝑞 ∇ 𝜍(𝐵(𝑥0 , 𝑟)) and 𝑔(𝑞) = 0 for 𝑞 ∈ 𝜍(𝐵(𝑥0 , 𝑠)). The function 𝑔 is called a bump
function around 𝑝 = 𝜍(𝑥0 ), because of the resemblance with speed bumps used to
reduce traffic.
BURALI-FORTI PARADOX: The Burali-Forti paradox demonstrates that the class of all
ordinals is not a set. If there were a set of all ordinals, 𝑂𝑟𝑑, then it would follow that 𝑂𝑟𝑑
was itself an ordinal, and therefore that 𝑂𝑟𝑑 ∇ 𝑂𝑟𝑑. Even if sets in general are allowed
to contain themselves, ordinals cannot since they are defined so that ∇ is well founded
over them.

BURNSIDE'S THEOREM: If 𝐺 is a finite group of order where 𝑝 and 𝑞 are prime

numbers, and 𝑎 and 𝑏 are non-negative integers, then 𝐺 is solvable. Hence each non-
Abelian finite simple group has order divisible by at least three distinct primes.

BUSEMANN FUNCTION: Given a ray, γ : 0, ∞ →X, the Busemann function is defined by

BUSEMANN’S THEOREM: Let 𝐾 be a convex body in 𝑛-dimensional Euclidean
space 𝑹𝒏 containing the origin in its interior. Let 𝑆 be an (𝑛 − 2)-dimensional linear
subspace of 𝑹𝒏 . For each unit vector 𝜃 in 𝑆 ⊥ , the orthogonal complement of 𝑆, let 𝑆𝜃
denote the (𝑛 − 1)-dimensional hyperplane containing 𝜃 and 𝑆. Define 𝑟(𝜃) to be the
(𝑛 − 1)-dimensional volume of 𝐾 ∩ 𝑆𝜃 . Let 𝐶 be the curve {𝜃𝑟(𝜃)} in 𝑆 ⊥ . Then 𝐶 forms
the boundary of a convex body in 𝑆 ⊥ .

BUTLER’S SPHERE THEOREM: Let a rigid sphere r = a be introduced into flow field of an
axis-symmetric irrotational flow in an incompressible inviscid fluid with no rigid
boundaries, characterized by the stream function Ψ0 = Ψ0 (r, θ)all of whose singularties
are at a distance greater than a from the origin, where Ψ0 = 0(r 2 ) at an origin, then the
stream function becomes

r a2
Ψ = Ψ0 − Ψ1 = Ψ0 r, θ − Ψ0 ,θ
a r

BUTTERFLY THEOREM: Let 𝑀 be the midpoint of a chord 𝑃𝑄 of a circle, through which

two other chords𝐴𝐵 and 𝐶𝐷 are drawn; 𝐴𝐷 and 𝐵𝐶 intersect chord 𝑃𝑄 at 𝑋 and 𝑌
correspondingly. Then 𝑀 is the midpoint of 𝑋𝑌.

𝑩𝑽 𝒂, 𝒃 SPACE: Let 𝑓 (𝑡) be complex valued function defined on the interval 𝑎, 𝑏 . We

take an arbitrary subdivision of [𝑎, 𝑏] as

𝑎 = 𝑡0 ≤ 𝑡1 ≤ 𝑡2 ≤ ⋯ ≤ 𝑡𝑛 = 𝑏,

and form the expression

 𝑛

𝑉𝑛 = 𝑓 𝑡𝑘 − 𝑓 𝑡𝑘−1 .
𝑘=1

It the aggregate of sums 𝑉𝑛 corresponding to all possible subdivisions of [𝑎, 𝑏] is

bounded, then the function 𝑓(𝑡) is defined as a function of bounded variation on the
interval 𝑎, 𝑏 , and the quantity

𝑉𝑎𝑏 𝑓 = 𝑠𝑢𝑝𝑘 𝑉𝑘,

is called the total variation of the function 𝑓 𝑡 .

We denote the space of all functions of bounded variation in the interval [𝑎, 𝑏] by
𝐵𝑉 𝑎, 𝑏 . That this is a linear space follows from the property that a function is of
bounded variation if and only if it can be expressed as the difference of two non-
decreasing functions.

C
𝑪𝟎 SPACE: This is the space of all sequences of complex numbers which converges to
zero, with the norm defined as in 𝑐. Since, if two sequences converge to zero, then their
linear combination also converges to zero, 𝑪𝟎 is a subspace of 𝑐. Thus 𝑐0 is a normed
linear space.

∞
For completeness, let 𝑓𝑛 be a fundamental sequence in 𝑐0 . Since 𝑐0 ⊂ 𝑐, 𝑓𝑛 𝑛=1 is also
a fundamental sequence is 𝑐, and has a limit 𝑐. By the theorem of double limits
𝑙𝑖𝑚𝑛→∞𝑓(𝑛) exists and is equal to 0. Hence

𝑓 = 𝑓(𝑛)∞
𝑛=1 lies in 𝑐.

It follows that 𝑐0 is complete normed linear space, i.e. a Banach space.

𝐂[𝟎, 𝟏] SPACE: The linear space 𝐶[0,1] of all continuous real function defined on the
closed unit interval [0,1] is a normed linear space with respect the norm of an element 𝑓
defined by

1
𝑓 = 𝑓 𝑡 𝑑𝑡,
0

But it is not Banach space.

𝐂[𝐚, 𝐛] SPACE: The space 𝐶 [𝑎, 𝑏] is the space of all bounded continuous functions
defined on the closed interval 𝑎, 𝑏 . We know that 𝐶 [𝑎, 𝑏] is a linear space with respect
to pointwise linear operations defined by

𝑓+𝑔 𝑥 = 𝑓 𝑥 + 𝑔 𝑥 and 𝛼𝑓 𝑥 = 𝛼𝑓 𝑥 ,

For all 𝑓, 𝑔 ∇ 𝐶 𝑎, 𝑏 , 𝑥 ∇ [𝑎, 𝑏] and all scalars 𝛼.

If 𝑓 is an arbitrary element of 𝐶 𝑎, 𝑏 , we define

𝑓 = 𝑠𝑢𝑝𝑥∇ 𝑎,𝑏 𝑓(𝑥)

Now for all 𝑓, 𝑔 ∇ 𝐶 𝑎, 𝑏 and all scalars 𝛼, we have:

i) 𝑓 = 𝑠𝑢𝑝𝑥∇ 𝑎,𝑏 𝑓(𝑥) , which is clearly non-negative.

ii) 𝑓 = 0 ⟺ 𝑠𝑢𝑝𝑥∇ 𝑎,𝑏 𝑓(𝑥) ,
⟺ 𝑓 𝑥 =, for each 𝑥 ∇ 𝑎, 𝑏 ,
⟺ 𝑓 = 0 (zero function)
iii) For any scalar 𝛼,
𝑎𝑓 = 𝑠𝑢𝑝𝑥∇ 𝑎,𝑏 𝑎𝑓 𝑥
= 𝑠𝑢𝑝𝑥∇ 𝑎,𝑏 𝑎 𝑓(𝑥)
= 𝛼 sup 𝑥 ∇ 𝑎, 𝑏 𝑓(𝑥)
= 𝛼 𝑓
iv) 𝑓 + 𝑔 = 𝑠𝑢𝑝𝑥∇ 𝑎,𝑏 𝑓+𝑔 𝑥
≤ 𝑠𝑢𝑝𝑥∇ 𝑎,𝑏 𝑓(𝑥) + 𝑔 𝑥
= 𝑠𝑢𝑝𝑥∇ 𝑎,𝑏 𝑓(𝑥) + 𝑠𝑢𝑝𝛼 ∇ 𝑎,𝑏 𝑔(𝑥)
= 𝑓 + 𝑔
cal: Calorie - a metric unit of energy. The equivalent SI unit of energy is joule. 1 cal
(lower case c) is the energy required to raise 1 gram of water by 1 °C at roughly 4.2
joules, while 1 Cal (capital C) is the energy required to raise 1 gram of water by 1 °C 4.2
kilojoules.

CALCULUS: Calculus is divided into two general areas: differential calculus and integral
calculus. The basic problem in differential calculus is to find the rate of change of a
function. Geometrically, this means finding the slope of the tangent line to a function at a
particular point; physically, this means finding the speed of an object if you are given its
position as a function of time. The slope of the tangent line to the curve 𝑦 = 𝑓 (𝑥) at a
𝑑𝑦
point (𝑥, 𝑓 (𝑥)) is called the derivative, written as 𝑦 ′ or . The reverse process of
𝑑𝑥

differentiation is integration (or anti-differentiation).

CALCULUS OF VARIATIONS: Calculus of variation is a development of calculus
concerned with problems in which a function is to be determined such that some
related definite integral achieves a maximum or minimum value. In calculus of
variations, the problem is to determine a curve 𝑦(𝑥) that minimizes (or maximizes) the
integral of a specified function over a specific range:
𝑏

𝐽= 𝑓 𝑥, 𝑦, 𝑦 ′ 𝑑𝑥
𝑎

To determine the function y, we will define a new quantity Y:

𝑌 = 𝑦 + 𝑒𝜌
where 𝑒 is a new variable, and 𝜌 can be any continuous function as long as it meets
these two conditions:
𝜌 𝑎 = 0, 𝜌 𝑏 =0
These conditions mean that the value for 𝑌 is the same as the value of 𝑦 at the two
endpoints of our interval 𝑎 and 𝑏. Then 𝐽 can be expressed as a function of 𝑒:
𝑏

𝐽(𝑒) = 𝑓 𝑥, 𝑌, 𝑌 ′ 𝑑𝑥
𝑎

To find the derivative:

𝑏
𝑑 𝑑
𝐽(𝑒) = 𝑓 𝑥, 𝑌, 𝑌 ′ 𝑑𝑥
𝑑𝑒 𝑑𝑒
𝑎
 𝑏
𝑑
= 𝑓 𝑥, 𝑌, 𝑌 ′ 𝑑𝑥
𝑑𝑒
𝑎
𝑏
𝜕𝑓 𝑑𝑥 𝜕𝑓 𝑑𝑌 𝜕𝑓 𝑑𝑌 ′
= + + 𝑑𝑥
𝜕𝑥 𝑑𝑒 𝜕𝑌 𝑑𝑒 𝜕𝑌 ′ 𝑑𝑒
𝑎

CANCELLATION: Any method of calculation the result of which, when compared to the
original form, attributes a number of components of the original form to mitigate the
effects of the rest (of those components), so that the result remains the same through
their omission. A common example invloves the numerator and denominator of
a fraction, and another invloves the omission of one logarithm symbol from each side of
the equation (when the each side consists of the logarithm only), even though the
former is a result of the equivalence of fractions and the latter the injective nature of
the logarithmic function. Failure of understanding these can result in the misuse of such
methods, and is commonplace amongst students.

CANCELLATION LAWS: Let ⟪ be a binary operation on a set 𝑆. The cancellation laws are
said to hold if, for all 𝑎, 𝑏 and 𝑐 in 𝑆,
(i) if 𝑎⟪𝑏 = 𝑎⟪𝑐, then b = c,
(ii) if 𝑏⟪𝑎 = 𝑐⟪𝑎, then b = c.
In a group G, the cancellation laws hold.
CANONICAL: Describes a representation of an object (e.g. an expression,
a transformation) in a way that is preferred, perhaps unique or considered natural due
to certain properties that it exhibits, even though there may be
other equivalent representations.

CANONICAL HEIGHT: The canonical height on an abelian variety is a height function that
is a distinguished quadratic form.
CANONICAL OR NORMAL FORM OF A REAL QUADRATIC FORMS: If X ′ AX is a real
quadratic form in 𝑛 variables, then there exists real non-singular linear transformation
𝑋 = PV which transforms X ′ AX to the form

𝑦12 + ⋯ + 𝑦𝑝2 − 𝑦𝑝+1

2
− ⋯ − 𝑦𝑟2
In the new form the given quadratic form has been expressed as a sum and difference
of the squares of new variables. This latter expression is called the canonical form or
normal form of the given quadratic form.

CANONICAL PRODUCT (COMPLEX ANALYSIS): Suppose 𝑓(𝑧) is an integral function of

finite order 𝜌 with an infinite number of zeroes 𝑧1 , 𝑧2 , 𝑧3 arranged in the order of
increasing modulus. Let 𝜌1 be the exponent of the convergence of zeros. We associate an
integer 𝜌 with sequence of zeros s.t.

𝜌1 , 𝜌1 ≠ 𝑖𝑛𝑡𝑒𝑔𝑒𝑟
∞

𝜌= 𝜌1 − 1, 𝜌1 𝑖𝑠 𝑖𝑛𝑡𝑒𝑔𝑒𝑟 𝑠. 𝑡. 𝑟𝑛 − 𝜌1 𝑖𝑠 𝑐𝑜𝑛𝑣𝑒𝑟𝑔𝑒𝑛𝑡
𝑛=1
𝜌1 , 𝑜𝑡𝑕𝑒𝑟𝑤𝑖𝑠𝑒

Where 𝜌1 denotes greatest integer less than or equal to 𝜌1

Evidently 𝜌1 − 1 ≤ 𝑝 ≤ 𝜌1 ≤ 𝜌 in any case.

Then the infinite product

∞
𝑧
𝐺 𝑧 = ,𝑝
𝑧𝑛
𝑛 =1

Converges uniformly and absolutely in any bounded closed domain of 𝑧 − 𝑝𝑙𝑎𝑛𝑒 which
does not contain any 𝑧𝑛 𝑛 = 1,2,3, … . Also 𝐺(𝑧) is an integral function and vanishes if
and only if 𝑧 is a zero of 𝑓 𝑧 . We call this product by the name Canocial product formed
with zero of 𝑓 𝑧 , the integer 𝜌 is called genus of it.

CANTILEVER: A beam or such similar structures which is anchored at only one end such
that it resists rotation under load.

CANTOR, GEORGE (1845–1918): George Cantor was a German Mathematician

responsible for the establishment of set theory and for profound developments in the
notion of the infinite. In 1873, he showed that the set of rational numbers is
denumerable. He also showed that the set of real numbers is not. Later he fully
developed his theory of infinite sets and so-called transfinite numbers. The latter part of
his life was clouded by repeated mental illness.
CANTOR INTERMEDIATE VALUE PROPERTY: A function 𝑓 ∶ 𝑅 → 𝑅 has the Cantor
intermediate value property if for every 𝑥, 𝑦 ∇ 𝑅 and for each perfect set 𝐾 between
𝑓 (𝑥) and 𝑓 (𝑦), there is a perfect set 𝐶 between 𝑥 and 𝑦 such that 𝑓 [𝐶] ⊂ 𝐾; the strong
Cantor intermediate value property if for every 𝑥, 𝑦 ∇ 𝑅 and for each perfect set 𝐾
between 𝑓 (𝑥) and 𝑓 (𝑦) there is a perfect set 𝐶 between 𝑥 and 𝑦 such that 𝑓 [𝐶] ⊂ 𝐾
and 𝑓I𝐶 is continuous; the weak Cantor intermediate value property if for every
𝑥, 𝑦 ∇ 𝑅 with 𝑓 (𝑥) < 𝑓 (𝑦) there exists a perfect set 𝐶 between 𝑥 and 𝑦 such that
𝑓 [𝐶] ⊂ (𝑓 (𝑥), 𝑓 (𝑦)); the perfect road if for every 𝑥 ∇ 𝑅 there exists a perfect set
𝑃 ⊂ 𝑅 having 𝑥 as a bilateral (i.e., two sided) limit point for which 𝑓I𝑃 is continuous at
𝑥.
CANTOR SET: Consider the closed interval [0, 1]. Remove the open interval that forms
1 2
the middle third, that is, the open interval , . From each of the remaining two
3 3

intervals, again remove the open interval that forms the middle third. The Cantor set is
the set that remains when this process is continued indefinitely. It consists of those real
numbers whose ternary representation 0. 𝑐1 𝑐2 𝑐3 … has each ternary digit 𝑐𝑖 equal to
either the digit 0 or 2.

CANTOR-BERNSTEIN-SCHROEDER THEOREM: The Cantor-Bernstein-Schroeder

theorem underlies the theory of transfinite cardinals. In an infinite set, there are subsets
of the exactly same cardinality. But then there are also different transfinite cardinalities.
So how does one compare infinite sets. Given two infinite sets 𝑨 and 𝑩, assume there is
a 1-1 correspondence between B and a subset of A. It is reasonable and viable to expect
that |𝑩| ≤ |𝑨|. Similarly, it is reasonable to say that |𝑩| < |𝐴| provided it is not true
that |𝑨| ≤ |𝑩|, which would hold if there was a 1-1 correspondence between A and a
subset of B. The theorem states as follows:

Let there be an injection 𝑓: 𝐴 → 𝐵 and another 𝑔: 𝐵 → 𝐴. Then there isa bijection

𝛼 ∶ 𝐴 → 𝐵. In other words, if |𝐴| ≤ |𝐵| and |𝐵| ≤ |𝐴|, then |𝐴| = |𝐵|.

CANTOR INTERSECTION THEOREM: It states that a decreasing nested sequence of non-

empty compact subsets of 𝑆 has nonempty intersection. In other words, supposing {𝐶𝑘 }
is a sequence of non-empty, closed and bounded sets satisfying

it follows that
CANTOR LEMMA: No interval [a, b] is countable.
CANTOR’S PARADOX: Suppose there exists an infinite set A containing the largest
possible number of elements. Cantor’s Diagonal Theorem shows that its power set has
more elements than A had. This proves there is no largest cardinal number.
CANTOR'S THEOREM: It states that, for any set A, the set of all subsets of 𝐴 (the power
set of 𝐴) has a strictly greater cardinality than 𝐴 itself.

CAPTURE-RECAPTURE SAMPLING: A method for estimating total population (usually of

animals) through 2 periods of capture (observation), assuming that the total number
has remained constant and that the probability of capture of any animal on any one visit
is constant and equal. In implementing this method, a researcher captures a number of
animals and mark them (or make sure that they can be identified in the future in some
way) and releases them, they come back to capture a number of animals and count the
proportion of animals that were captured both times. By assuming that the proportion
of animals recaptured in an unbiased estimate of the proportion of animals marked in
the first capture, combine this with the number of animals marked on the first visit, we
can easily calculate the estimated number of the total population. In practice, it is
difficult to ensure that the animal population remains constant (or roughly the same)
while ensuring that the chance of recapture is truly equal amongst the population since
the former requires the visits to be relatively close to each other so that the population
does not change significantly and yet the latter condition requires the visits to be
sufficiently space apart so that the locations of the animals are truly randomized.

CARATHÉODORY MEASURABLE SET: Given an outer measure µ∗ , we define 𝐸 ⊆ 𝑋 to be

Carathéodory measurable if 𝜇 ∗ 𝐴 = 𝜇 ∗ 𝐴 ∩ 𝐸 + 𝜇 ∗ 𝐴\𝐸 for any 𝐴 ⊆ 𝑋. As µ∗ is sub-
additive, this is equivalent to 𝜇 ∗ 𝐴 ≥ 𝜇 ∗ 𝐴 ∩ 𝐸 + 𝜇 ∗ 𝐴\𝐸 as the other inequality is
automatic. Measurability by Lebesgue and Carathéodory are equivalent.

CARATHÉODORY'S THEOREM: It states that if 𝑈 is a simply connected open subset of the

complex plane C, whose boundary is a Jordan curve 𝛤 then the Riemann map 𝑓: 𝑈 → 𝐷
from 𝑈 to the unit disk 𝐷 extends continuously to the boundary, giving a
homeomorphism 𝐹 : 𝛤 → 𝑆 1 from 𝛤 to the unit circle 𝑆 1 .
CARD: It is the abbreviation for cardinality of a set. (Card(X) is also written
#𝑋, ♯𝑋 𝑜𝑟 |𝑋|).
CARDINALITY OF A SET: For a finite set 𝐴, the cardinality of 𝐴, denoted by 𝑛(𝐴), is the
number of elements in 𝐴. The notation #(𝐴) or |𝐴| is also used. For subsets 𝐴, 𝐵 and 𝐶 of
some universal set 𝐸,
 𝑛(𝐴 ∪ 𝐵) = 𝑛(𝐴) + 𝑛(𝐵) – 𝑛(𝐴 ∩ 𝐵)
 𝑛(𝐴 ∪ 𝐵 ∪ 𝐶) = 𝑛(𝐴) + 𝑛(𝐵) + 𝑛(𝐶) – 𝑛(𝐴 ∩ 𝐵) – 𝑛(𝐴 ∩ 𝐶) – 𝑛(𝐵 ∩
𝐶) + 𝑛(𝐴 ∩ 𝐵 ∩ 𝐶).
CARDINAL NUMBER: A number that gives the number of elements in a set. If two sets
can be put in one-to-one correspondence with one another they have the same cardinal
number or cardinality. For finite sets, the cardinal numbers are 0, 1, 2, 3, …, but infinite
sets require new symbols to describe their cardinality like aleph and aleph-null.
CARDIOID: A cardioids is the curve traced out by a point on the circumference of a circle
rolling round another circle of the same radius. Its equation, in which 𝑎 is the radius of
each circle, may be taken in polar coordinates as 𝑟 = 2𝑎(1 + 𝑐𝑜𝑠𝜃)(– 𝜋 < 𝜃 ≤ 𝜋). In
the figure, OA = 4a and OB = 2a.

CARMICHAEL'S THEOREM: It states that for 𝑛 greater than 12, the 𝑛th Fibonacci
number 𝐹(𝑛) has at least one prime divisor that does not divide any earlier Fibonacci
number.

The only exceptions for 𝑛 up to 12 are:

 𝐹(1) = 1 and 𝐹(2) = 1, which have no prime divisors

 𝐹(6) = 8 whose only prime divisor is 2 (which is 𝐹(3))
 𝐹(12) = 144 whose only prime divisors are 2 (which is 𝐹(3)) and 3 (which is
𝐹(4))
If a prime 𝑝 is a divisor of 𝐹(𝑛) that does not divide any 𝐹(𝑚) with 𝑚 < 𝑛, then 𝑝 is
called a characteristic factor or a primitive prime divisor of 𝐹(𝑛).
Carmichael's theorem says that every Fibonacci number, apart from the exceptions
listed above, has at least one primitive prime divisor.

CARTAN: It is extended Einstein's General relativity to Einstein-Cartan theory, using

Riemannian-Cartan geometry instead of Riemannian geometry. This extension provides
affine torsion, which allows for non-symmetric curvature tensors and the incorporation
of spin-orbit coupling.
CARTAN–HADAMARD THEOREM: It states that a connected, simply connected complete
Riemannian manifold with non-positive sectional curvature is diffeomorphic to Rn via
the exponential map; for metric spaces, the statement that a connected, simply
connected complete geodesic metric space with non-positive curvature in the sense of
Alexandrov is a globally CAT(0) space.

CARTESIAN COORDINATE SYSTEM: The rectilinear coordinate system used for any
(finite and natural) number of dimensions. The number of axes used is commonly seen
as a basic way of "defining" the number of dimensions.

These axes must all be mutually perpendicular and the components (the coordinates) of
a point are the shortest distances of a point to the hyperplane consisting the axes
(except the axis related to the component in question). Thus, the x-coordinate of a point
is the shortest distance from the point to any point on the line (plane) consisting of
the y-axis (both y and z-axes) in 2D (3D).

While rectilinear coordinate system is a more descriptive name of its function, it is more
commonly called Cartesian, named after Cartesius, the Latin name of René Descartes.
CARTESIAN PRODUCT: The Cartesian product 𝐴 × 𝐵, of sets 𝐴 and 𝐵, is the set of all
ordered pairs (𝑎, 𝑏), where 𝑎 ∇ 𝐴 and 𝑏 ∇ 𝐵. Similarly, the Cartesian product
𝐴 × 𝐵 × 𝐶 of sets 𝐴, 𝐵 and 𝐶 can be defined as the set of all ordered triples (𝑎, 𝑏, 𝑐),
where 𝑎 ∇ 𝐴, 𝑏 ∇ 𝐵 and 𝑐 ∇ 𝐶.

CARTESIAN TENSOR: Cartesian tensors are extensively used in a variety of branches of

continuum mechanics, such as fluid mechanics and elasticity. In classical continuum
mechanics, the space of concentration is usually 3-dimensional Euclidean space, as is
the tangent space at each point. If we put a ceiling on the local coordinates to be
Cartesian coordinates with the same scale centered at the point of concentration, the
metric tensor is the Kronecker delta. This means that there is no requirement to
distinguish covariant and contravariant components, and furthermore there is no need
to distinguish tensors and tensor densities. All Cartesian-tensor indices are written as
subscripts. Cartesian tensors achieve considerable computational simplification at the
cost of generality and of some theoretical insight.

CASORATI-WEIERSTRASS THEOREM: If 𝑓 ∶ 𝐷(𝑃, 𝑟0 ) \ {𝑃} → 𝐶 is holomorphic and 𝑃 is

an essential singularity of 𝑓, then 𝑓(𝐷(𝑃, 𝑟) \ {𝑃}) is dense in 𝐶 for any 0 < 𝑟 < 𝑟0 .
CASTING OUT NINES: It is a method of verification of integer arithmetic by checking that
the answer matches the equivalent calculations under modulo arithmetic.
The modulus of 9 is chosen simply because of our decimal systems (base 10).
The difference of 1 between the modulus (9) and the base (10) allows for an easy
conversion of a number to modulo 9. (By adding the constituent digits of a number
represented in base 10.)

CATEGORICAL DATA: Observations that indicate categories to which individuals belong

rather than measurements or values of variables. Often such data consists of a summary
that shows the numbers or frequencies of individuals in each category. This form of data
is known as a 'frequency table' or, if there are two or more categorized features, as 'a
contingency table'.
CATEGORICAL VARIABLE: A random variable with values which are categories. (e.g.
nationality, gender etc.)

CATEGORY OF A SET: A topological characterization of the "massiveness" of a set. 𝑨

subset 𝑬 of a topological space 𝑿 is said to be of the first category in 𝑿 if it can be
expressed as a finite or countable union of nowhere dense sets in 𝑿, otherwise 𝑬 is said
to be of the second category. This terminology is, however, not universal: some authors
use the name second category for complements in 𝑿 of sets of the first category. In the
case of a Baire space, a more appropriate name for such sets is residual. A non-empty
closed set of real numbers, in particular an interval, is not of the first category in itself.
This result generalizes to any complete metric space, it is called Baire category theorem
and has wide application in analysis. The role of a set of the first category in topology is
analogous to that of a null set in measure theory. However, in 𝑹 a set of the first
category can be a set of full (Lebesgue) measure, while there are (Lebesgue) null sets
which are residual.
CATENARY: The curve in which an ideal flexible heavy rope or chain of uniform density
hangs between two points. The equation of the curve is 𝑦 = 𝑐 𝑐𝑜𝑠𝑕 (𝑥/𝑐).

A catenary is a curve represented by the formula

1
𝑦= 𝑎 𝑒 𝑥/𝑎 + 𝑒 −𝑥/𝑎
2
The value of 𝑎 is the 𝑦 intercept. The catenary can also be represented by the hyperbolic
cosine function 𝑦 = 𝑐𝑜𝑠𝑕 𝑥.
CATENOID: Catenoid is the surface generated by rotating a catenary around its axis of
symmetry.
CAUCHY, AUGUSTIN-LOUIS (1789–1857): Augustin-Louis Cauchy was one of the most
important mathematicians of the early nineteenth century and a dominating figure in
French mathematics. His work ranged over vast areas of mathematics, in almost 800
papers, but he is chiefly remembered as one of the founders of rigorous mathematical
analysis. Using the definition of limit as it is now known; he developed sound definitions
of continuity and convergence. He was also a pioneer in the theory of functions of a
complex variable.

CAUCHY CONVERGENCE CONDITION: 1. For sequences - a necessary and sufficient

condition for the convergence of a sequence: an infinite sequence is convergent, if and
only if, for any positive number 𝜀, there is always a number 𝑁(𝜀) (𝑁 is a function of 𝜀),
such that the difference between any two terms after the 𝑁th term is less than 𝜀.

2. For series - a necessary and sufficient condition for the convergence of a series: An
infinite series 𝛴𝑎𝑛 is convergent, if and only if, for any positive number 𝜀 , there is an
integer 𝑁(𝜀) (𝑁 is a function of 𝜀 ), such that

𝑎𝑛+1 + 𝑎𝑛+2 + − − − + 𝑎𝑝 <∇

for all integers 𝑛 > 𝑁 and 𝑝 is a positive integer.

CAUCHY DISTRIBUTION: A symmetric continuous probability distribution with infinite

support.

CAUCHY ESTIMATE: If 𝑓 is a holomorphic on a region containing the closed disc 𝐷(𝑃, 𝑟)

𝜕𝑘 𝑀 · 𝑘!
and if |𝑓| ≤ 𝑀 on 𝐷 (𝑃, 𝑟), then 𝑓(𝑃) ≤ .
𝜕𝑥 𝑘 𝑟𝑘

CAUCHY HADAMARD THEOREM: Consider the formal power series in one complex
variable z of the form

Where 𝑎, 𝑐𝑛 ∇ ℂ. Then the radius of convergence of ƒ at the point 𝑎 is given by

where lim sup denotes the limit superior, the limit as 𝑛 approaches infinity of
the supremum of the sequence values after the 𝑛th position. If the sequence values are
unbounded so that the lim sup is ∞, then the power series does not converge near 𝑎,
while if the lim sup is 0 then the radius of convergence is ∞, meaning that the series
converges on the entire plane.

CAYLEY HAMILTON THEOREM: If 𝛼 ∶ 𝐶 2 → 𝐶 2 is a linear map, let us write 𝑄(𝑡) =

𝑑𝑒𝑡(𝑡𝜄 − 𝛼). Then we have 𝑄(𝑡) = 𝑡 2 + 𝑎𝑡 + 𝑏 where 𝑎, 𝑏 ∇ 𝐶. The Cayley–
Hamilton theorem states that 𝛼 2 + 𝑎𝛼 + 𝑏𝜄 = 𝑂 or, more briefly that Q α = O.

CAUCHY INTEGRAL FORMULA FOR THE DERIVATIVE OF AN ANALYTIC FUNCTION: If a

function 𝑓(𝑧) is analytic within and on a closed contour 𝐶 𝑎𝑛𝑑 𝑎 is any point lying in it,
then

1 𝑓 𝑧 𝑑𝑧
𝑓′ 𝑎 = ∫𝐶
2𝜋𝑖 (𝑧 − 𝑎)²

𝑡
CAUCHY MACLAURIN’S INTEGRAL TEST: Let 𝐹 𝑡 = ∫𝑎 𝑓 𝑥 𝑑𝑥 for 𝑎 ≤ 𝑡 < ∞. If
∞
lim𝑡→∞ 𝐹(𝑡) exists and is equal to 𝑙 ∇ 𝑅, the improper integral ∫𝑎 𝑓 𝑥 𝑑𝑥 is convergent
to 𝑙 otherwise, it is divergent.
CAUCHY MEAN VALUE THEOREM: If real valued functions 𝑓(𝑥) and 𝑔(𝑥) are such that
 𝑓 𝑥 , 𝑔(𝑥) are continuous in the closed interval [𝑎, 𝑏]
 𝑓 𝑥 , 𝑔(𝑥) are differentiable in the open interval (𝑎, 𝑏)
𝑓 ′ (𝑐) 𝑓 𝑏 −𝑓(𝑎)
then there exist at least one value of 𝑥 = 𝑐 ∇ (𝑎, 𝑏) such that 𝑔 ′ (𝑐) = 𝑔 .
𝑏 −𝑔(𝑎)

CAUCHY RATIO TEST: Also known simply as the ratio test. It is a method of deciding
the convergence of a series through the general ratio of one term of the series to the
next.

CAUCHY RIEMANN EQUATIONS: For an analytic function 𝑓(𝑧) = 𝑢 + 𝑖𝑣 of the complex

variable 𝑧 = 𝑥 + 𝑖𝑦 the Cauchy–Riemann equations linking the real and imaginary
parts of the function are
𝜕𝑢 𝜕𝑣 𝜕𝑣 𝜕𝑢
= 𝑎𝑛𝑑 =−
𝜕𝑥 𝜕𝑦 𝜕𝑥 𝜕𝑦
CAUCHY SCHWARZ–BUNYAKOVSKII INEQUALITY: For vectors 𝑥 and 𝑦 in an inner
product space 𝑉 let us define ||𝑥|| = ⟨ 𝑥, 𝑥 ⟩𝑎𝑛𝑑 ||𝑦|| = ⟨ 𝑦, 𝑦 ⟩ then we have
⌌𝑥, 𝑦⌍ ≤ 𝑥 𝑦 with equality if and only if 𝑥 and 𝑦 are scalar multiple each other.

CAUCHY SCHWARZ INEQUALITY FOR INTEGRALS: If 𝑓(𝑥), 𝑔(𝑥) are real functions then
{∫ [𝑓(𝑥)𝑔(𝑥)]𝑑𝑥}2 ≤ {∫ [𝑓(𝑥)]2 𝑑𝑥}{∫ [𝑔(𝑥)]2 𝑑𝑥}
if all these integrals exist.
CAUCHY SCHWARZ INEQUALITY FOR SUMS: If 𝑎𝑖 and 𝑏𝑖 are real numbers, 𝑖 = 1, 2, … , 𝑛
then

𝑛 𝑛 𝑛

𝑎𝑖 𝑏𝑖 ≤ 𝑎𝑖2 𝑏𝑖2
𝑖=1 𝑖=1 𝑖=1

CAUCHY SEQUENCE: A sequence ⌌𝑎𝑛 ⌍ for which the metric 𝑑(𝑎𝑛 , 𝑎𝑚 ), where 𝑚 > 𝑛,
satisfies lim𝑛 →∞ 𝑑 𝑎𝑛 , 𝑎𝑚 = 0.
Cauchy sequences converge when they are defined on the set of real numbers, but do
not necessarily converge on the set of rational numbers.
CAUCHY’S CONDENSATION TEST: If the function 𝑓(𝑥) is positive for all positive integral
values of 𝑛 and continually decreases as 𝑛 increases, then the two infinite series 𝑓(𝑛)
and 𝑎𝑛 𝑓(𝑎𝑛 ) are either both convergent or both divergent, 𝑎 being a positive integer
greater than 1.
CAUCHY’S FIRST THEOREM ON LIMITS: If ⌌𝑥𝑛 ⌍ is a sequence of positive reals such that
lim𝑛→∞ 𝑥𝑛 = 𝑙, then
𝑥1 + 𝑥2 + 𝑥3 + − − − + 𝑥𝑛
lim =𝑙
𝑛 →∞ 𝑛
CAUCHY’S GENERAL PRINCIPLE OF CONVERGENCE: A sequence ⌌𝑥𝑛 ⌍ of real numbers
converges if and only if it is a Cauchy sequence.
CAUCHY’S INEQUALITY COMPLEX ANALYSIS): If 𝑓(𝑧) is analytic within and on a circle
𝑛 𝑀𝑛 !
𝐶, given by 𝑧 − 𝑎 = 𝑅 𝑎𝑛𝑑 𝑖𝑓 𝑓(𝑧) ≤ 𝑀 for every 𝑧 𝑜𝑛 𝐶, then 𝑓 (𝑎) ≤ 𝑅𝑛

CAUCHY’S INTEGRAL THEOREM: For a closed curve C and an analytic function 𝑓(𝑧),

𝑓 𝑧 𝑑𝑧 = 0

CAUCHY’S LEMMA: If 𝐺 is a finite group and 𝑝 is a prime number that divides the order
of 𝐺, then 𝐺 must contain an element of order 𝑝.
 1
CAUCHY’S ROOT TEST: Let 𝑢𝑛 be a series of positive terms such that lim𝑛→∞ 𝑢𝑛 𝑛 = 𝑙.
Then
1. 𝑢𝑛 converges if 𝑙 < 1.
2. 𝑢𝑛 diverges if 𝑙 > 1.
3. The test fails if 𝑙 = 1.
CAUCHY’S SECOND THEOREM ON LIMITS: If ⌌𝑥𝑛 ⌍ is a sequence of positive reals such that
lim𝑛→∞ 𝑥𝑛 = 𝑙, then
𝑛
lim 𝑥1 𝑥2 𝑥3 − − − 𝑥𝑛 = 𝑙
𝑛→∞

CAUCHY’S THEOREM: If 𝐺 is a finite group and 𝑝 is a prime number which is a divisor of

the order of the group𝐺, then 𝐺 contains an element of order 𝑝. This implies that there
must be a subgroup of 𝐺 whose order is 𝑝.
CAUCHY’S THEOREM COMPLEX ANALYSIS): If a function 𝑓(𝑧) is analytic and single
valued inside and on a simple closed contour 𝐶, then ∫𝑐 𝑓 𝑧 𝑑𝑧 = 0

CAUSE VARIABLE: Also known as an independent variable, an explanatory variable, a

predictor or predictor variable or regressor.

CAUSTIC: The envelope of rays refracted or reflected by a geometric object.

CAVALIERI’S PRINCIPLE: Two geometrical figures whose cross sections are the same as
each other, at the same distance away from some reference line/lines (plane/planes)
have the same area (volume). As an example, this explains why triangles whose bases
have the same length, and have the same height, have also the same area, regardless of
its shape.

CAYLEY, ARTHUR (1821–95): Arthur Cayley was a British mathematician who

contributed greatly to the resurgence of pure mathematics in Britain in the nineteenth
century. He published over 900 papers on many aspects of geometry and algebra. He
conceived and developed the theory of matrices, and was one of the first to study
abstract groups.
CAYLEY–HAMILTON THEOREM: The characteristic polynomial 𝑝(𝜆) of an 𝑛 × 𝑛 matrix
𝑨 is defined by 𝑝(𝜆) = 𝑑𝑒𝑡(𝑨 – 𝜆𝑰). The following result about the characteristic
polynomial is called the Cayley–Hamilton Theorem:
Theorem: If the characteristic polynomial 𝑝(𝜆) of an 𝑛 × 𝑛 matrix 𝑨 is 𝑝(𝜆) =

(– 1)𝑛 (𝜆𝑛 + 𝑏𝑛–1 𝜆𝑛–1 + ··· + 𝑏1 𝜆 + 𝑏0 ),

then 𝐴𝑛 + 𝑏𝑛–1 𝐴𝑛–1 + ··· + 𝑏1 𝐴 + 𝑏0 𝑰 = 𝟎.

CAYLEY-HAMILTON THEOREM FOR RINGS: Let 𝑈 be be an 𝑛 × 𝑛 matrix with
coefficients in a unital commutative ring 𝑅 and let 𝜒𝑈 (𝑡) be the characteristic
polynomial of 𝑈. Then 𝜒𝑈 (𝑈) = 0.
CAYLEY REPRESENTATION THEOREM: Every group is isomorphic to a group of
permutations.
cdf: It is the abbreviation for cumulative distribution function. Short form of
Cumulative distribution function of a probability distribution.

CEILING FUNCTION: It is the function on real numbers whose value is always rounded
up, if the argument is not already an integer. The function leaves integers unchanged.

CELESTIAL MECHANICS: The study of motions of celestial bodies such as stars, planets
and comets etc.

CELL: Categories of data divided into rectangular arrays through more than
one variable. It is essentially a conventional use for the analogue of groups in 1
dimension.

CELSUIS: Represented by the symbol °C. It is based on dividing the difference between
the freezing point and boiling point of water into 100 equal "degrees.

CENSORED OBSERVATIONS: In statistics, these are the observations that are made
incomplete systematically due to the nature of the procedure (possible) for observation
or the objects under study.

CENSUS: A survey of the entire population.

CENTESIMAL MEASURE: A metric angular measure dividing a right angle into 100
centesimal degrees.

CENTI-: An SI prefix for one one hundredth.

CENTRAL ANGLE: An angle between two radii (of a circle or sphere).

CENTRAL AXIS: The axis of symmetry of a right-circular cone

CENTRAL CONIC: A central conic is a conic with a centre of symmetry, and thus an
ellipse or a hyperbola. The conic with equation
𝑎𝑥 2 + 2𝑕𝑥𝑦 + 𝑏𝑦 2 + 2𝑔𝑥 + 2𝑓𝑦 + 𝑐 = 0
is central if and only if
𝑎𝑏 ≠ 𝑕2 .
CENTRAL DIFFERENCE: If {(𝑥𝑖 , 𝑓𝑖 )}, 𝑖 = 0, 1, 2, … is a given set of function values with
𝑥𝑖+1 = 𝑥𝑖 + 𝑕, 𝑓𝑖 = 𝑓(𝑥𝑖 ) then the central difference at 𝑓𝑖 is defined by
𝑓𝑖+1 – 𝑓𝑖−1 𝑓 𝑥𝑖+1 − 𝑓(𝑥𝑖−1 )
=
2 2
CENTRAL DIFFERENCE APPROXIMATION: The most common numerical approximation
to the derivative of a function 𝑓(𝑥) is to take the gradient of the chord joining the point
and another point where 𝑥 has been increased by a small amount 𝑕. The central
difference approximation uses the chord joining the two points whose 𝑥 values are a
small amount 𝑕 from the value 𝑥0 , giving
𝑓 𝑥0 + 𝑕 − 𝑓 𝑥0 − 𝑕
𝑓 ′ (𝑥0 ) ≈
2𝑕
CENTRAL FIELD: The family of curves 𝑦 = 𝑦 𝑥, 𝑐 is said to form a central field over a
domain D if together they cover the whole D without intersecting each other.
CENTRALIZER: Let 𝐺 be a group. The centralizer 𝐶(𝑕) of an element 𝑕 of 𝐺 is the
subgroup of 𝐺 defined by 𝐶(𝑕) = {𝑔 ∇ 𝐺 ∶ 𝑔𝑕 = 𝑕𝑔}.
CENTRAL LIMIT THEOREM: If independent samples of size n are taken from a
population the distribution of the sample means (known as the sampling distribution of
the mean) will be approximately normal for large n. The mean of the sampling
distribution is the population mean and its variance is the population variance divided
by n. Note that there is no need for the distribution in the population sampled to be
normal. The theorem is important because it allows statements to be made about
population parameters by comparing the results of a single experiment with those that
would be expected according to some null hypothesis.

CENTRAL TENDENCY: A common measure in summary statistics bsed around the loose
idea that we can assign one location to represent the locations of a number of objects
considered as one. Thus, there is not just one but rather a number of slightly different
concepts which fits the description of central tendency. It is what is commonly referred
to by the similarly loose idea of an average.

CENTRE (CENTRE OF SYMMETRY): A point about which a geometric figure is in some

way self-similar. It can either be a centre of rotational symmetry, or the intersections of
multiple lines/planes of reflective symmetry, or even the intersections of medians in
a triangle.

CENTRE OF A FUZZY SET: if the mean value of all points at which the membership
function of the fuzzy set achieves its maximum value is finite, then define this mean
value as the center of the fuzzy set; if the mean value equals positive (negative) infinite,
then the center is defined as the smallest (largest) among all points that achieve the
maximum membership value.

CENTRE OF A GROUP: The set of elements which are commutative with every elements
of the group. Note that there must be at least one element in the centre, the identity.
Whereas the centre can be as large as the group, such a group is called an Abelian group.

CENTRE OF BUOYANCY: The centre of gravity of the body of water that an object
displaces.

CENTRE OF CURVATURE: Given a curve, the centre of curvature of a point on this curve
is the centre of a circle which "locally" (for a neighborhood of that point) describes the
curvature of the curve.

CENTRE OF GRAVITY: A point through gravity can be considered to be acting, instead of

individually on the point masses or acting on the body as a whole. It is a different
concept from the centre of mass, although due to the small differences in gravitational
strength in most context, they can be considered close approximates of each other.

CENTRE OF MASS: The point (not necessarily within the object) which is the
weighted average of the point masses of a body (or the set of infinitely many point
masses through integration), which can be used to calculate linear motion of a rigid
body as if all of the object's mass are at that point (the centre of mass) only.
(Considering the object as a particle.) It is also known as the barycentre. A point q ∇ M is
called the center of mass of the points if it is a point of global minimum
of the function

Such a point is unique if all distances are less than radius of convexity.
For uniform bodies:
2
TRIANGULAR LAMINA: along median from vertex
3
𝑟 𝑠𝑖𝑛𝛼
CIRCULAR ARC, RADIUS r, ANGLE AT CENTRE 2α: from centre
α
2𝑟 𝑠𝑖𝑛𝛼
SECTOR OF CIRCLE, RADIUS r, ANGLE AT CENTRE 2𝛼: from centre
3α

CENTRE OF ROTATION: The invariant point in a rotation.

CENTRIPETAL COMPONENT: A component of an object's acceleration corresponding to

a centripetal force.

CENTRIPETAL FORCE: A force perpendicular to the velocity of an object which causes

the object to travel on a curved (not straight) path.

CENTROID: The center of mass of an object. For a one-dimensional uniform object of

length L, the centroid is the midpoint of the line segment. For a triangle, the centroid is
the intersecting point of its three medians. The centroid of a symmetrical figure is the
center of symmetry. For any other irregular shaped two-dimensional object, the
centroid is the point where this single support can balance this object. Generally, the
centroid of a two- or three-dimensional object is found by using double or triple
integrals.
CESARO’S THEOREM: If ⌌𝑥𝑛 ⌍ and ⌌𝑦𝑛 ⌍ are two sequences of positive reals such that
lim𝑛→∞ 𝑥𝑛 = 𝑙 and lim𝑛→∞ 𝑦𝑛 = 𝑙 ′ , then
𝑥1 𝑦𝑛 + 𝑥2 𝑦𝑛−1 + 𝑥3 𝑦𝑛−2 + − − − + 𝑥𝑛 𝑦1
lim = 𝑙𝑙 ′
𝑛 →∞ 𝑛

CEVIAN: Any line segment joining a vertex of a triangle to a point on the infinite line
containing the opposite side.
c.g.s. UNITS: A system of units based on centimeters, grams and seconds. (Instead
of meters, kilograms and hours for the SI units). It has been largely superseded by the
use of SI units nowadays.

CHAIN: Let 𝐵 ⊆ 𝐴, where 𝐴 is ordered by ≤. 𝐵 is a chain in 𝐴 if any two elements of 𝐵

are comparable. That is, 𝐵 is a linearly ordered subset of 𝐴.
CHAIN COMPLEXES: A chain complex 𝐶 ∗ is a (doubly infinite) sequence (𝐶𝑖 ∶ 𝑖 ∇ 𝑍) of
modules over some unital ring, together with homomorphisms 𝜕𝑖 ∶ 𝐶𝑖 → 𝐶𝑖−1 for each
𝑖 ∇ 𝑍, such that 𝜕𝑖 ∘ 𝜕𝑖 + 1 = 0 for all integers 𝑖.
CHAIN MAP: Let 𝐶 ∗ and 𝐷 ∗ be chain complexes. A chain map 𝑓: 𝐶 ∗ → 𝐷∗ is a sequence
𝑓𝑖 ∶ 𝐶𝑖 → 𝐷𝑖 of homomorphisms which satisfy the commutativity condition
𝜕𝑖 ∘ 𝑓𝑖 = 𝑓𝑖−1 ∘ 𝜕𝑖 for all 𝑖 ∇ 𝑍.
CHAIN RULE: The following rule that gives the derivative of the composition of two
functions: If 𝑕(𝑥) = 𝑓(𝑔(𝑥)) for all x, then 𝑕′(𝑥) = 𝑓′(𝑔(𝑥))𝑔′(𝑥).
char: It is the abbreviation for characteristic of a ring.
CHARACTERS: A character 𝜒 of an Abelian group 𝐺 is a function which assigns to each
𝑎 ∇ 𝐺 a complex number 𝜒(𝑎) of absolute value 1 and satisfies 𝜒(𝑎𝑏) = 𝜒(𝑎)𝜒(𝑏)for all
𝑎, 𝑏 ∇ 𝐺. The product 𝜒 = 𝜒1 𝜒2 of two characters 𝜒1 𝑎𝑛𝑑 𝜒2 is delïned by 𝜒(𝑎) =
𝜒1 (𝑎)𝜒2 (𝑏), and 𝜒 is also a character of 𝐺. Thus all the characters of 𝐺 form an Abelian
group 𝐶(𝐺), which is called the character group of 𝐺. The identity element of the
character group is the identity character (or principal character) 𝜒 such that 𝜒(𝑎) = 1
for all 𝑎 ∇ 𝐺. If 𝐺 is finite, then 𝐺 ≈ 𝐶(𝐺). This implies the duality 𝐺 = 𝐶[𝐶(𝐺)].
CHARACTERIZATION OF POLYNOMIALS: The order of a zero of a polynomial equals the
order of its first non-vanishing derivative.

CHARACTERISTIC FUNCTION: Suppose A is a subset of a set X. Then the function

1, 𝑖𝑓 𝑥 ∇ 𝐴
𝜒𝐴 𝑥 =
0, 𝑖𝑓 𝑥 ∈ 𝐴
is the characteristic function for 𝐴.
CHARACTERISTIC OF A FIELD: The smallest positive whole number 𝑛 such that the sum
of the multiplicative identity added to itself 𝑛 times equals the additive identity. If no
such 𝑛 exists, the field is said to have characteristic zero.
CHARACTERISTIC OF A RING: Let 𝑅 be a ring, and let 𝑟 ∇ 𝑅. We may define 𝑛. 𝑟 for all
natural numbers 𝑛 by recursion on 𝑛 so that 1. 𝑟 = 𝑟 and 𝑛. 𝑟 = (𝑛 − 1). 𝑟 + 𝑟 for all
𝑛 > 0. We define also 0. 𝑟 = 0 𝑎𝑛𝑑 (−𝑛). 𝑟 = −(𝑛. 𝑟) for all natural numbers n. Then
(𝑚 + 𝑛). 𝑟 = 𝑚. 𝑟 + 𝑛. 𝑟, 𝑛. (𝑟 + 𝑠) = 𝑛. 𝑟 + 𝑛. 𝑠, (𝑚𝑛). 𝑟 = 𝑚. (𝑛. 𝑟), (𝑚. 𝑟)(𝑛. 𝑠) =
(𝑚𝑛). (𝑟𝑠) for all integers 𝑚 an 𝑛 and for all elements 𝑟 and 𝑠 of 𝑅. In particular,
suppose that 𝑅 is a unital ring. Then the set of all integers n satisfying 𝑛. 1 = 0 is an
ideal of 𝑍. Therefore there exists a unique nonnegative integer 𝑝 such that 𝑝𝑍 = {𝑛 ∇
𝑍 ∶ 𝑛. 1 = 0} . This integer 𝑝 is referred to as the characteristic of the ring 𝑅, and is
denoted by 𝑐𝑕𝑎𝑟 𝑅.
CHARACTERISTIC POLYNOMIAL: Let 𝑨 be a square matrix. Then 𝑑𝑒𝑡 (𝑨 – 𝜆𝑰) is a
polynomial in 𝜆 and is called the characteristic polynomial of 𝑨. The equation
𝑑𝑒𝑡(𝑨 – 𝜆𝑰) = 0 is the characteristic equation of 𝑨, and its roots are the characteristic
values of 𝑨.
CHARACTERISTIC SUBSPACES OF A MATRIX: Suppose λ is an eigenvalue of a square
matrix 𝐴. Then every non-zero vector 𝑋 satisfying the equation

A − λI = 0 ………. 1

is an eigenvector of 𝐴 corresponding to the eigenvalue λ. If the matrix 𝐴 − λI is of rank 𝑟,

then the equation (1) will possess 𝑛 − 𝑟 linarly independent solutions. Each non-zero
linear combination of these solutions is also a solution of (1) and therefore it will be an
eigenvector of 𝐴. The set of all these linear combinations is a subspace of 𝑉𝑛 provided we
add zero vector also to this set. This subspace of 𝑉𝑛 is called characteristic subspace of 𝐴
corresponding to the eigenvalue λ. It is nothing but the column null space of the matrix
𝐴 − λI. Its dimension 𝑛 − 𝑟 is the geometric multiplicity of the eigenvalue λ.

CHARACTERISTIC VALUE: Let 𝑨 be a square matrix. The roots of the characteristic

equation 𝑑𝑒𝑡(𝑨 – 𝜆𝑰) = 0 are called the characteristic values of 𝑨. Then 𝜆 is a
characteristic value of 𝑨 if and only if there is a non-zero vector 𝒙 such that 𝑨𝒙 = 𝜆𝒙.
Let A = aij be any 𝑛-rowed square matrix and λ an indeterminate. The matrix
n×n

𝐴 − λI is called the characteristic matrix of 𝐴 where 𝐼 is the unit matrix of order 𝑛.

a11 − λ a12 … a1n

a21 a22 − λ … a2n
Also the determinant A − λI = which is an ordinary
… … … …
an1 an2 … ann − λ
polynomial of λ of degree 𝑛, is called the characteristic polynomial of 𝐴. The equation
A − λI = 0 is called the characteristic equation of A and the roots of this equation are
called the characteristic roots or characteristics values or eigen values or latent roots or
proper values of the matrix A. The set of the eigen values of A is called the spectrum of
A.

If λ is a characteristic roots of the matrix 𝐴, then A − λI = 0 and the matrix A − λI is

singular. Therefore there exists a non-zero vector 𝑋 such that

A − λI X = 0 𝑜𝑟 𝐴𝑋 = λX.

CHARACTERISTIC VECTOR: Any vector 𝒙 such that 𝑨𝒙 = 𝜆𝒙 is called a characteristic

vector corresponding to the characteristic value 𝜆. If λ is a characteristic root of an 𝑛 × 𝑛
matrix 𝐴, then a non-zero vector 𝑋 such that 𝐴𝑋 = λX is called a character vector or
eigenvector of 𝐴 corresponding to the characteristic root λ. λ is a characteristic root of a
matrix 𝐴 if and only if there exists a non-zero vector 𝑋 such that 𝐴𝑋 = λX.
If 𝑋 is a characteristic vector of a matrix 𝐴 corresponding to the characteristic value λ,
then 𝑘𝑋 is also a characteristic vector of 𝐴 corresponding to the same characteristic
value λ. Here 𝑘 is any non-zero scalar.

If 𝑋 is a characteristic vector of a matrix 𝐴, then 𝑋 cannot corresponding to more than

one characteristic value A.

The characteristic vectors corresponding to distinct characteristic roots of a matrix are

linear independent.

If λ1 be a characteristic root of order 𝑡 of the characteristic equation A − λI = 0, then 𝑡

is called the algebraic multiplicity of λ1 . If s be the number of linearly independent
eigenvector corresponding to the eigenvalue λ1 , then s is called the geometric
multiplicity of λ1 . In this case number of linearly independent solutions of A − λ1 I X =
0 will be s and the matrix A − λ1 I will be rank n − s.

The geometric multiplicity of a characteristic root cannot exceed its algebraic

multiplicity i.e., s ≤ t.

The characteristic roots of a Hermitian matrix are real.

The characteristic roots of a skew-Hermitian matrix are either pure imaginary or zero.
The characteristic roots of a real symmetric matrix are either pure imaginary or zero,
for every such matrix is skew- Hermitian.

The characteristic roots of a unitary matrix are of unit modulus.

The characteristic roots of an orthogonal matrix are of unit modulus.

CHARACTER OF THE GROUP: If 𝐺 is a Hausdorff locally compact Abelian group we say

that a continuous group homomorphism 𝜒 ∶ 𝐺 → 𝑆 1 is a character of the group. We
write 𝐺 for the set of such characters and 𝑕𝑥, 𝜒𝑖 = 𝜒(𝑥) for all 𝑥 ∇ 𝐺 and 𝜒 ∇ 𝐺 .

CHART: Let 𝑆 be a surface in 𝑅 3 . A chart on 𝑆 is an injective regular parametrized

surface 𝜍: 𝑈 → 𝑅 3 with image 𝜍(𝑈) ⊂ 𝑆. A collection of charts 𝜍𝑖 : 𝑈𝑖 → 𝑅 3 on 𝑆 is said
to cover 𝑆 if 𝑆 = ⋃𝑖 𝜍𝑖 𝑈𝑖 . In that case the collection is called an atlas of 𝑆. A coordinate
map, a coordinate chart, or simply a chart, of a manifold is an invertible map between a
subset of the manifold and a simple space such that both the map and its inverse
preserve the desired structure. For a topological manifold, the simple space is
some Euclidean space Rn and interest focuses on the topological structure. This
structure is preserved by homeomorphisms, invertible maps that are continuous in
both directions. In the case of a differentiable manifold, a set of charts called
an atlas allows us to do calculus on manifolds. Polar coordinates, for example, form a
chart for the plane R2 minus the positive x-axis and the origin.

CHEBYSHEV APPROXIMATION: Let 𝐷 be a bounded closed subset of the complex plane,

and 𝑓(𝑧) a continuous function on 𝐷. Then there exists a polynomial 𝜋𝑛 𝑧 of degree 𝑛
such that 𝑚𝑎𝑥𝑧𝜀𝐷 𝑓 𝑧 − 𝜋𝑛 𝑧 attains the infimum 𝐸𝑛 𝑓 . The polynomial 𝜋𝑛 𝑧 is
unique and is called the best approximation polynomial (in the sense of Chebyshev). If
𝐷 is simply connected and 𝑓 𝑧 is single-valued and holomorphic on 𝐷, then 𝜋𝑛 𝑧
converges to 𝑓 𝑧 uniformly on 𝐷. Moreover, in this case there exist a number 𝑀 that
does not depend on 𝑛 and a number 𝑅 > 1 such that

𝑓 𝑧 − 𝜋𝑛 𝑧 ≤ 𝑀 𝑅𝑛 .

CHEBYSHEV, PAFNUTY LVOVICH (1821–94): Pafnuty Lvovich Chebyshev was a Russian

mathematician and founder of a notable school of mathematicians in St Petersburg. His
name is remembered in results in algebra, analysis and probability theory. In number
theory, he proved that, for all 𝑛 > 3, there is at least one prime between 𝑛 and 2𝑛 – 2.
CHEBYSHEV’S DIFFERENTIAL EQUATION: The differential equation

d2 y dy
1 − 𝑥2 − x + n2 y = 0
dx 2 dx

is called Chebyshev’s differential equation.

CHEBYSHEV’S INEQUALITIES (PROBABILITY): Chebyshev proved a number of

inequalities relating to the maximum proportion of distributions which could lie beyond
1
a certain point: 𝑃 𝑋 − 𝜇𝑥 > 𝑘𝜍 ≤ 𝑘 2 says that the probability a random variable 𝑋 lies

more than 𝑘 standard deviations from its mean is not more than 1/𝑘 2 . If X is a random
𝐸{𝑔 𝑋 }
variable and 𝑔(𝑋) is always ≥ 0 then 𝑃 𝑔(𝑋) ≥ 𝑘 ≤ says that for a non-negative
𝑘

function of a random variable, the probability the function takes a value at least 𝑘 can be
no more than the mean of the function divided by 𝑘.
While these inequalities are very weak statements in that most distributions do not
come close to the limit specified, it is very useful sometimes to be able to identify an
upper limit that it is impossible for a probability to exceed.
CHEBYSHEV’S INEQUALITY (MEASURE THEORY): If 𝑓 is non-negative & summable, then
1
𝜇 𝑥 ∇ 𝑋: 𝑓 𝑥 > 𝑐 < 𝑐 ∫ 𝑓𝑑𝜇. Let (𝑓𝑛 ) be monotonically ↑ sequence of µ −summable

functions on 𝑋. Define 𝑓(𝑥) = lim𝑛 →∞ 𝑓𝑛 (𝑥) (allowing the value +∞).

1. If all integrals ∫ 𝑓𝑛 𝑑µ are bounded, then 𝑓 is summable and ∫ 𝑓𝑛 𝑑µ =

lim𝑛→∞ ∫ 𝑓𝑛 𝑑µ .

2. If lim𝑛→∞ ∫ 𝑓𝑛 𝑑µ = ∞ then function 𝑓 is not summable.

Let 𝑓 be a measurable non-negative function attaining only finite values. 𝑓 is summable

if and only if 𝑠𝑢𝑝∫ 𝑓 𝑑µ < ∞, where the supremum is taken over all finite-measure set 𝐴
such that 𝑓 is bounded on 𝐴.

CHEBYSHEV’S POLYNOMIALS: The Chebyshev polynomials of first kind, Tn (x) and

second kind Un (x) are defined by

Tn x = cos(ncos −1 x)

and Un x = sin(ncos −1 x)

where n is a non-negative integer.

Sometimes the Chebyshev polynomials of the second kind is defined by

1
Un x = sin (n + 1)cos−1 x) / [ 1 − x 2 ] Un+1 (x)
[ 1 − x2

Chebyshev’s polynomials are also known as Teheibehef, Tbbicheff or Tshebysheff’s

polynomials.

CHEBYSHEV'S SUM INEQUALITY: It states that if

and

then

Similarly, if

and

then

CHEBYSHEV’S THEOREM (NUMBER THEORY): For any positive integer greater than 𝑛
there is always a prime between 𝑛 and 2𝑛.
CHEBYSHEV’S THEOREM (STATISTICS): For a random variable, whatever the
distribution, with 𝐸(𝑋) = 𝜇, 𝑉𝑎𝑟(𝑋) = 𝜍 2 the proportion of values which lie within 𝑘
1
standard deviations of the mean will be at least 1 − 𝑘 2 .

CHINESE POSTMAN PROBLEM: The problem of finding the least weighted circuit in a
connected graph. If an Eulerian cycle exists than such a cycle is the solution to the
problem, otherwise, it is necessary to repeat at least one edge and the problem can be
more complicated.
CHINESE REMAINDER THEOREM: Suppose 𝑛1 , . . . , 𝑛𝑘 are positive integers that
are pairwise coprime. Then, for any given sequence of integers 𝑎1 , . . . , 𝑎𝑘 , there exists
an integer 𝑥 solving the following system of simultaneous congruences.

Furthermore, all solutions 𝑥 of this system are congruent modulo the product,
𝑁 = 𝑛1 , . . . , 𝑛𝑘 . Hence

CHI-SQUARED DISTRIBUTION: Chi-squared distribution is a type of non-negative

continuous probability distribution, normally written as the χ2-distribution, with one
parameter υ called the degrees of freedom. The distribution is skewed to the right and
has the property that the sum of independent random variables each having a χ2-
distribution also has a χ2-distribution. It is used in the chi-squared test for measuring
goodness of fit, in tests on variance and in testing for independence in contingency
tables. It has mean υ and variance 2υ.
CHI-SQUARED TEST: A test, normally written as the χ2-test, to determine how well a set
of observations fits a particular discrete distribution or some other given null
hypothesis. The observed frequencies in different groups are denoted by 𝑂𝑖 , and the
expected frequencies from the statistical model are denoted by 𝐸𝑖 . For each 𝑖, the value
(𝑂𝑖 – 𝐸𝑖 )2 /𝐸𝑖 is calculated, and these are summed. The result is compared with a chi-
squared distribution with an appropriate number of degrees of freedom. The number of
degrees of freedom depends on the number of groups and the number of parameters
being estimated. The test requires that the observations are independent and that the
sample size and expected frequencies exceed minimum numbers depending on the
number of groups. Chi-Square is a squared standard normal variable that has zero mean
and unit variance but is more often used for the sum of v independent squared standard
normal variables. The mean of this quantity is v, and this is termed the "degrees of
freedom". Common uses include:
 To test a sample variance against a known value;
  To test the homogeneity of a set of sample means where the population variance
is known;
 To test the homogeneity of a set of proportions;
 To test the goodness-of-_t of a theoretical distribution to an observed frequency
distribution;
 To test for association between categorized variables in a contingency table;
 To construct a confidence interval for a sample variance
CHRISTOFFEL’S EXPANSION: P′n = 2n − 1 Pn−1 + 2n − 5 Pn−3 + 2n − 9 Pn−5 + ⋯
The last terms of the series being 3P1 or P0 according as n is, even or odd.

CHRISTOFFEL’S FIRST SUMMATION FORMULA: Christoffel’s first summation formula is

given as

n P n +1 x P n y −P n +1 y P n x
r=0 2r + 1 Pr x Pr y = (n + 1) .
(x−y)

CHRISTOFFEL’S SECOND SUMMATION FORMULA: Christoffel’s second summation

formula is given as

𝑛
1 + n + 1 [Pn+1 x Qn y − Pn x Qn+1 y
(2𝑟 + 1) Pr x Qr y =
y−x
𝑟=0

CHROMATIC NUMBER: For a graph or map 𝐺, the maximum number of colours needed
so that all regions touching one another (meeting at an edge or a vertex) are in a
different colour is the chromatic number, denoted by 𝜒(𝐺). The Four Colour Theorem
proved that 𝜒(𝐺) ≤ 4 for all planar graphs.
CIRCLE: The conic produced by slicing a right-circular cone at right angles to its central
axis
CIRCLE OF CONVERGENCE: A circle in the complex plane with the property that
∞
𝑖=1 𝑎𝑖 (𝑧 − 𝑧0 )𝑖 converges for all 𝑧 within a distance 𝑅 > 0 of 𝑧0 and diverges for all 𝑧
for which |𝑧 − 𝑧0 | > 𝑅 . 𝑅 is the radius of convergence and if the power series
converges for the whole of the complex plane then 𝑅 is infinite. The case where 𝑅 = 0
is trivial since convergence only occurs when 𝑧 − 𝑧0 = 0. For points on the
circumference of the circle the series may either converge or diverge.
CIRCLE OF CURVATURE: A circle which describes correctly the curvature of a curve for
the neighborhood of a point on the curve.

CIRCLE THEOREM: Let 𝑓(𝑧) be the complex potential of motion of a two-dimensional

irrotational flow of an incompressible, inviscid fluid with no rigid boundaries and 𝑓(𝑧)
has no singularities within the circle z = a. If a circular cylinder, classified by its cross-
section the circle z = a, be introduced into the field of flow, the complex potential
becomes

𝑤 = 𝑓 𝑧 + 𝑓 (𝑎2 /𝑧)

where f is the complex conjugate form of f.

Note that the angular velocity (rate of rotation) at any point of the flowing fluid is equal
to the half the curl of the velocity at the point. The vector curl q is termed the vector or
the vorticity of the fluid at a point.

Body forces are those which act equally on all the matter within a small element of
volume and the total force is proportional to the size of the volume element or mass of
the fluid considered. The simplest example of such a force is the gravitational force.
Centrifugal force is treated as a body force when the coordinate system is rotating, but it
is not treated as a body force if the coordinates system is stationary. Thus the body can
be specified as per unit mass of the fluid. It may vary from point to point of the fluid, and
at any point it may have different values at different instances of times. The body forces
have three components and can be represented by a vector function of position and
time.

Surface forces are viewed s acting on any surface divided in the fluid including its
boundaries. Each point in the fluid may belong to several surfaces. Several such surf e
forces seem to co-exist at the point.

Surface forces have a direct molecular origin, decrease rapidly with increase of distance
between interacting elements. These are considerable only when that distance is of the
separation of molecules of the fluid. They are negligible unless there is direct
mechanical contact between the interacting elements. These are the forces exerted by
adjacent portions of the fluid upon one another and are in the nature of action and
reaction. The specification of these force are obtained from the concept of stress tensor
which has nine components.

CIRCUIT: A closed path on a graph.

CIRCULAR: Of or related to a circle.

CIRCULAR CONE: A cone whose base is a circle. It is normally known as cone.

CIRCULAR CYLINDER: A cylinder whose base is a circle. It is normally known as
cylinder.

CIRCULAR DATA: The class of cyclic directional data in statistics.

CIRCULAR FUNCTIONS: Another name for the trigonometric functions. Trigonometric

functions are based on parameterization of a circle just as hyperbolic functions are
based on parameterization of hyperbolae. This usage highlights the analogue between
the two.

CIRCULAR HELIX: A helix which lies on the surface of a circular cylinder is called a
circular helix or right circular helix.

CIRCULAR MEASURE: Also known as angular measure.

CIRCULAR MOTION: The motion of an object whose path forms a circle.

CIRCULAR SECTIONS: When the intersection of a plane and a quadric is a circle, the
intersection is called a circular section. In general, circular sections are cut off by two
systems of parallel planes through a quadric. The point of contact on the tangent plane
parallel to these is an umbilical point of the quadric.

CIRCUMCENTRE: The centre of a circle which goes through all vertices of a polygon.

CIRCUMCIRCLE: A circle whose circumference contains all points of the polygon.

CIRCUMFERENCE: The length of the closed curve of a circle.

CIRCUMSCRIBED: It is the act of enclosing a geometric figure with a
(minimal) circle or sphere.

CIS: A commonly defined function in the study of complex numbers. 𝑐𝑖𝑠(𝜃) = 𝑐𝑜𝑠(𝜃) +
𝑖𝑠𝑖𝑛(𝜃) Incidentally, 𝑐𝑖𝑠(𝜃) = 𝑒𝑥𝑝(𝑖𝜃). As such, 𝑐𝑖𝑠 is seldom used apart from a
particular stage of learning in school as most favour the use of the simpler 𝑒𝑥𝑝(𝑖𝜃).

CLAIRAUT'S EQUATION: A family of differential equations of the form

which can be solved by differentiating the whole equation.

CLAIRAUT’S THEOREM DIFFERENTIAL GEOMETRY): If a geodesic on a surface of

revolution cuts the meridian through any point P on it at an angle Ψ, then 𝑢 sin Ψ is
constant where 𝑢 the distance of the point P from the axis is.

CLAPEYRON FORMULA: Let us assume that the energy density function

𝑊 ℯ1, ℯ2 , … … . . 𝑒6 can be expanded in a power series

2 𝑊= 𝑐0 + 2𝑐𝑖 𝑒𝑖 , + 𝑐𝑖𝑗 𝑒𝑖 𝑒𝑗 + … … .,

Neglecting constant terms 𝑐𝑜 and third and higher order terms in the strains.
Differentiating partially (1) w.r.t 𝑒𝑖 , we have

2𝜕𝑊
= 2𝑐𝑖 + 𝑐𝑖𝑗 𝑒𝑗 + 𝑐𝑗𝑖 𝑒𝑖
𝜕𝑒𝑖

2𝜕𝑊 1
Or = 𝜏𝑖 = 𝑐𝑖 + 2 𝑐𝑖𝑗 + 𝑐𝑗𝑖 𝑒𝑗 ∵ 𝑖 𝑖𝑠 𝑑𝑢𝑚𝑚𝑦
𝜕𝑒 𝑖

If 𝜏𝑖 = 0 𝑤𝑖𝑡𝑕 𝑠𝑡𝑟𝑎𝑖𝑛 𝑒𝑗 , 𝑤𝑒 𝑕𝑎𝑣𝑒 𝑐𝑖 = 0 and thus we get,

1
𝑊 = 2 𝑐𝑖𝑗 𝑒𝑗

𝜕𝑊 1
𝐻𝑒𝑛𝑐𝑒 𝜏𝑖 = = 𝑐 + 𝑐𝑗𝑖 𝑒𝑗 .
𝜕𝑒𝑖 2 𝑖𝑗
Also the coefficient in the generalized Hooke’s law is symmetric of the strain- energy
density function W, with the property stated above, exists.

Thus we can write 𝜏𝑖 = 𝑐𝑖𝑗 𝑒𝑗 , where 𝑐𝑖𝑗 = 𝑐𝑗𝑖

This gives

1
𝑊 = 2 𝜏𝑖𝑒 𝑖 𝑖 = 1,2, … . .6 . This formula is known as the Clapeyron’s formula.

CLASS: A collection of objects, not necessary a set. The distinction generally is that
certain operations on sets are not allowed so as to allow us to refer to a (usually large)
collection of objects without causing contradictions such as Russell's paradox.

CLASS FREQUENCY: It is the number of occurrence in a class.

CLASSICAL MECHANICS: A loosely defined term generally considered to be the study

of motion of objects before the drastically different quantum mechanics. Thus classical
mechanics include Newtonian, Lagrangian, Hamiltonian and arguably, relativistic
mechanics.

CLASSIFICATION THEOREM: Every finite simple group is isomorphic to one of the

following groups:

 A cyclic group with prime order;

 An alternating group of degree at least 5;
 A simple group of Lie type, including both

 the classical Lie groups, namely the simple groups related to

the projective special linear, unitary, symplectic or orthogonal
transformations over a finite field;
 the exceptional and twisted groups of Lie type.

 The 26 sporadic simple groups.

The classification theorem has applications in many branches of mathematics, as

questions about the structure of finite groups (and their action on other mathematical
objects) can sometimes be reduced to questions about finite simple groups.
CLASS INTERVALS: The intervals in which data fall into a particular class.

CLASS OF A SURFACE: If the equation of the surface be 𝑥 = 𝑓 𝑢, 𝑣 ; 𝑦 = 𝑔 𝑢, 𝑣 ; 𝑧 =

𝑢, 𝑣 ; 𝑧 = 𝑕(𝑢, 𝑣) then the surface is said to be of class 𝑟 if the functions 𝑓, 𝑔, 𝑕 are
single valued as well continuous and possess partial derivatives of rth order.
CLOSED (GRAPH THEORY): A walk, trail or path which finishes at its starting point is
closed.
CLOSED BALL: If (𝑀, 𝑑) is a metric space, a closed ball is a set of the form 𝐷(𝑥; 𝑟) : =
{𝑦 ∇ 𝑀 : 𝑑(𝑥, 𝑦) ≤ 𝑟}, where 𝑥 is in 𝑀 and 𝑟 is a positive real number, the radius of
the ball. A closed ball of radius 𝑟 is a closed 𝒓-ball. Every closed ball is a closed set in the
topology induced on 𝑀 by 𝑑. Note that the closed ball 𝐷(𝑥; 𝑟) might not be equal to
the closure of the open ball 𝐵(𝑥; 𝑟).
CLOSED CURVE: A continuous plane curve that has no ends or, in other words, that
begins and ends at the same point. In other words, a closed curve is a curve that
completely encloses an area.
CLOSED FIGURE: It is a shape or a curve that begins and ends at the same point.
CLOSED GRAPH: Let 𝑋 and 𝑌 be metric spaces. A correspondence 𝛤 ∶ 𝑋 → 𝑌 is said to
be closed at 𝑥 ∇ 𝑋, if, for any sequence (𝑥𝑚 ) ∇ 𝑋 and (𝑦𝑚 ) ∇ 𝑌, 𝑥𝑚 → 𝑥, 𝑦𝑚 ∇ 𝛤(𝑥𝑚 )
(for each 𝑚) and 𝑦𝑚 → 𝑦 imply 𝑦 ∇ 𝛤(𝑥). 𝛤 is said to have a closed graph if it is closed
at every 𝑥 ∇ 𝑋.

CLOSED GRAPH THEOREM: Let 𝐵 𝑎𝑛𝑑 𝐵′ be Banach spaces and let 𝑇 be a linear
treansformation of 𝐵 into 𝐵 ′ . Then 𝑇 is a continuous mapping of and only if its graph is
closed.

CLOSED INTERVAL: The closed interval [𝑎, 𝑏] is the set of all real numbers between 𝑎
and 𝑏 including 𝑎 and 𝑏 i.e. {𝑥|𝑥 ∇ 𝑹, 𝑎 ≤ 𝑥 ≤ 𝑏}.

CLOSED RANGE THEOREM: Let 𝑋 and 𝑌 be Banach spaces, 𝑇: 𝐷(𝑇) → 𝑌 a closed linear
operator whose domain 𝐷(𝑇) is dense in 𝑋, and 𝑇 ′ the transpose of 𝑇. The theorem
asserts that the following conditions are equivalent:

 𝑅(𝑇), the range of 𝑇, is closed in 𝑌.

 𝑅(𝑇 ′ ), the range of 𝑇 ′ , is closed in 𝑋 ′ , the dual of 𝑋.
  .

 .

CLOSED REGION: It is a set of points containing all limit points and is connected.

CLOSED SET: The complement of an open set in a metric space or in a topological space
is called a closed set. Therefore, a subset of a topological space is called closed if it is the
complement of an open set

CLOSED SURFACE: A closed surface is a surface that completely encloses a volume of

space. For example, a sphere is a closed surface, but a teacup is not.
CLOSED WALK: Let (𝑉, 𝐸) be a graph. A walk 𝑣0 𝑣1 𝑣2 . . . 𝑣𝑛 in the graph is said to be
closed if 𝑣0 = 𝑣𝑛 . Thus a walk in a graph is closed if and only if it starts and ends at the
same vertex.
CLOSURE OF A SET: The closure of an open set A is obtained by including in it all limit
points of the set A.
CLOSURE PROPERTY: An arithmetic operation obeys the closure property with respect
to a given set of numbers if the result of performing that operation on two numbers
from that set will always be a member of that same set. For example, the operation of
addition is closed with respect to the integers, but the operation of division is not.
Operation Natural Numbers Integers Rational Numbers Real Numbers
Addition Closed Closed Closed Closed
Subtraction Not Closed Closed Closed Closed
Division Not Closed Not Closed Closed Closed
Root extraction Not Closed Not Closed Not Closed Not Closed
CLUSTER POINT: 𝑥 ∗ ∇ 𝑅 is a cluster point of the sequence ⌌𝑥𝑛 ⌍ iff, given any 𝜀 > 0, and
any positive integer 𝑚, there exists 𝑛 > 𝑚 such that |𝑥𝑛 − 𝑥 ∗ | < 𝜀. The real number 𝑥 ∗
is a cluster point of the sequence ⌌𝑥𝑛 ⌍ iff there exists a subsequence ⌌𝑥𝑛 𝑖 ⌍ of ⌌𝑥𝑛 ⌍, such
that 𝑥𝑛 𝑖 → 𝑥 ∗ . If the sequence ⌌𝑥𝑛 ⌍ converges to 𝑥 ∗ ∇ R, then 𝑥 ∗ is a cluster point; in fact
the only cluster point of ⌌𝑥𝑛 ⌍. CM: An abbreviation for the centre of mass of an object.

CLUSTER SAMPLING: Where a population is geographically scattered it is reasonable to

divide it into regions from which a sample is taken, and then a sample of individuals is
taken from those regions only. The result is that the individuals in the final sample
appear as clusters in the original population, but the costs of taking the sample are
much lower than doing a full random sampling process. There are different strategies
possible at both stages of the sampling process.

COAXIAL: The property of having the same axis.

CODIMENSION: The codimension of a submanifold is the dimension of the ambient

space minus the dimension of the submanifold.
CODING THEORY: The area of mathematics concerned with the encryption of messages
to ensure security during transmission, and with the recovery of information from
corrupted data. With increasing use of the internet and other electronic
communications to conduct business, this is one of the developing areas of mathematics
research, for example encryption using numbers based on the product of very large
primes.

CO-DOMAIN: The co-domain of a function is a superset of the image of the function. It is

sometimes known as the range of the function. The co-domain serves as a constraint as
to what values a function can take, thus there can be elements of the co-domain which is
not the value of the function but not vice versa.

COEFFICIENT: The quantity with which we multiply the variable in question.

coefficient matrix: A matrix formed from a system of linear equations using

the coefficients of the equations only (omitting the variables). Also known as
an augmented matrix
COEFFICIENT OF CORRELATION: A name for several related methods which measure
the relationship between two sets of data.

COEFFICIENT OF DETERMINATION: A measure of how much of the variation in the

data can be accounted for by the statistical model, for the purpose of inferring the likely
level of determination of outcomes.

COEFFICIENT OF FRICTION: A dimensionless quantity, the ratio of the friction force to

the normal reaction as determined by a number of other factors.

COEFFICIENT OF KURTOSIS: A measure of how much the data concentrates around

the mean. Informally, it measured the "pointiness"/"peakedness" of the
distribution curve.

COEFFICIENT OF RESTITUTION: A dimensionless quantity, the ratio of the speed of

separation to the speed of approach of 2 objects as determined by other factors. (e.g. the
material the objects consist of and their arrangement.)

COEFFICIENT OF SKEWNESS: A measure of the asymmetry of a distribution.

COEFFICIENT OF VARIATION: coefficient of variation is a measure of dispersion equal to

the standard deviation of a sample divided by the mean. The value is a dimensionless
quantity, not dependent on the units or scale in which the observations are made, and is
often expressed as a percentage.
i+j
COFACTORS: The minor Mij multiplied by −1 is called the cofactor of the element
aij . We shall denote the cofactor of an element by the corresponding capital letter. With
this notation, cofactor of aij = Aij = (−1)i+j Mij . For example,

The cofactor of the element a21 = A21 = (−1)2+1 M21

a12 a13
=− a a33 ,
32

a11 a13
The cofactor of the element a32 = A32 = − a a23 ,
21

a22 a23
The cofactor of the element a11 = A11 = − a a33 , and so no
32
Thus the cofactor of any element aij = (−1)i+j × the determinant obtained by leaving
the row and the column passing through that element.

In terms of the notation of cofactors, we have

∆= a11 A11 + a12 A12 + a13 A13 ,

or ∆= a21 A11 + a22 A22 + a23 A23 ,

or ∆= a13 A13 + a23 A23 + a33 A33 , and so on.

Therefore, in a determinant the sum of the products of the elements of any row or
column with the corresponding cofactors is equal to the value of the determinant.

COFUNCTION IDENTITIES: A set of trigonometric identites which relates trigonometric

functions to their cofunctions, e.g. sin and cos, sec and csc, tan and cot.

COINCIDENT: The property of two geometric figures to have all points in common.

CO-INTIAL VECTOR: The vectors which have the same initial point are called co-initial
vectors.

COLATITUDE: The angle between the vector and the polar axis (z-axis) in a
Spherical polar coordinate system, whereas latitude is the value subtracting colatitude
from 90 degrees, the smallest angle between the radius vector and the plane
perpendicular to the polaw axis through the origin. (the equator)

COLLINEAR: A set of three or more points is collinear if they all lie on the same line.
COLLINEARITY OF THREE POINTS: The necessary and sufficient condition for three
points with position vectors 𝒂, 𝒃, 𝒄 to be collinear is that there exist three scalars x, y, z
not all zero, such that x𝒂 + y𝒃+ z𝒄=0 where x + y + z= 0.

COLLINEAR OR PARALLEL VECTORS: Vectors having the same line of action or having
the lines of action parallel to one another are called collinear or parallel vectors.

COLLISION: The interaction of two objects with each other through contact transitioned
from a state of non-contact prior.

COLLUSION: Several customers may collaborate and only one of them may stand in the
queue.

COLUMN RANK OF A MATRIX: Let A = aij be any 𝑚 × 𝑛 matrix. Each of the 𝑛 rows of 𝐴
consists of 𝑚 elements. Therefore the column vectors of 𝐴 are 𝑚-vectors. These column
vectors of 𝐴 will span a subspace 𝐶 to 𝑉𝑚 . This subspace 𝐶 is called the column space of
the matrix 𝐴. The dimension 𝑐 of 𝐶 is called the column rank of 𝐴. In other words the
column rank of a matrix 𝐴 is equal to the maximum number of linearly independent
column of 𝐴.

COLUMN SPACE: Let 𝜏 ∇ 𝐿(𝐹 𝑚 , 𝐹 𝑛 ) be represented by a matrix 𝑇. So the columns of 𝑇

are 𝜏 (𝑒1 ), . . . , 𝜏 (𝑒𝑚 ). Thus the image of 𝜏 is the subspace of 𝐹 𝑛 spanned by the columns
of 𝑇. We call this the column space of 𝑇 and define the column rank of 𝑇 to be the
dimension of this column space.
COMBINATION: A selection of all or part of a set of objects, where the order of the
selection is not significant.
COMBINATORICS: The branch of mathematics that includes the study of permutations
and combinations. The area of mathematics concerned with counting strategies to
calculate the ways in which objects can be arranged to satisfy given conditions.
COMMON DIFFERENCE: A constant that is added to each term to produce an arithmetic
sequence. It is the constant difference between any two successive terms in an
arithmetic sequence.

COMMON LOGARITHM: A base-10 logarithm.

COMMON RATIO: It is a constant that is multiplied to each term to produce a geometric
sequence. It is the ratio of successive terms in a geometric sequence.
COMMON REFINEMENT OF TWO PARTITIONS: Let 𝑃 and 𝑃 ∗ be two partitions of a closed
interval [𝑎, 𝑏]. Then 𝑃∗∗ is said to be a common refinement of 𝑃 and 𝑃∗ if 𝑃∗∗ = 𝑃 ∪ 𝑃 ∗ .
COMMON TANGENT: One single line that forms a tangent to two or more different
curves. The term is also used for the distance between two tangential points.
COMMUNICATION CHANNEL: A communication channel is comprised of an alphabet
𝐴 = {𝑎1 , . . . , 𝑎𝑞 } and a set of forward channel probabilities of the form
𝑃[𝑎𝑗 𝑟𝑒𝑐𝑒𝑖𝑣𝑒𝑑 | 𝑎𝑖 𝑤𝑎𝑠 𝑠𝑒𝑛𝑡] such that for every 𝑖:
𝑞

𝑃[𝑎𝑗 𝑟𝑒𝑐𝑒𝑖𝑣𝑒𝑑 | 𝑎𝑖 𝑤𝑎𝑠 𝑠𝑒𝑛𝑡] = 1

𝑗 =1

A communication channel is memoryless if for all vectors 𝑥 = 𝑥1 . . . 𝑥𝑛 and

𝑐 = 𝑐1 . . . 𝑐𝑛 it holds that
𝑛

𝑃[𝑥 𝑟𝑒𝑐𝑒𝑖𝑣𝑒𝑑 | 𝑐 𝑤𝑎𝑠 𝑠𝑒𝑛𝑡] = 𝑃[𝑥𝑖 𝑟𝑒𝑐𝑒𝑖𝑣𝑒𝑑 | 𝑐𝑖 𝑤𝑎𝑠 𝑠𝑒𝑛𝑡]

𝑖=1

COMMUTATIVE: The binary operation ⟪on a set 𝑆 is said to be commutative

if,
𝑎⟪ 𝑏 = 𝑏 ⟪𝑎 ∀ 𝑎 , 𝑏 ∇ 𝑆
COMMUTATIVE ALGEBRA: An algebra A is called commutative algebra if

𝑓𝑔 = 𝑔𝑓, 𝑓𝑜𝑟 𝑎𝑙𝑙 𝑓, 𝑔 𝜖𝐴

COMMUTATIVE BINARY OPERATIONS: A binary operation ∗ on a set 𝐴 is said to be

commutative if 𝑥 ∗ 𝑦 = 𝑦 ∗ 𝑥 for all elements 𝑥 and 𝑦 of 𝐴. Example The operations of
addition and multiplication on the set R of real numbers are commutative, since
𝑥 + 𝑦 = 𝑦 + 𝑥 and 𝑥 × 𝑦 = 𝑦 × 𝑥 for all real numbers 𝑥 and 𝑦. However the
operation of subtraction is not commutative, since 𝑥 − 𝑦 ≠ 𝑦 − 𝑥 in general. (Indeed
the identity 𝑥 − 𝑦 = 𝑦 − 𝑥 holds only when 𝑥 = 𝑦.)
COMMUTATOR: The element 𝑥, 𝑦 = 𝑥 −1 𝑦 −1 𝑥𝑦 for 𝑥, 𝑦 in the group. It has the
property that [𝑥, 𝑦] and [𝑦, 𝑥] must be commutative, because one is the inverse of the
other, and therefore both products will be equal to the identity.
COMMUTE: Let ⟪ be a binary operation on a set 𝑆. The elements 𝑎 and 𝑏 of 𝑆 are said to
commute if 𝑎⟪ 𝑏 = 𝑏 ⟪𝑎.
COMPACTIFICATION: A compactification of a space 𝑋 is a compact space 𝑌 such that 𝑋 is
homeomorphic to a dense subspace of 𝑌. A compactification of the Euclidean interval
(0, 1) is the Euclidean interval [0, 1].

COMPACT OPEN TOPOLOGY: Let 𝑋 and 𝑌 be two topological spaces. Let 𝐶(𝑋, 𝑌) be the
set of all continuous functions from 𝑋 to 𝑌. The topology on generated by all sets of the
form 𝑆(𝐴, 𝑈) = { 𝑓 ∇ 𝐶(𝑋, 𝑌) | 𝑓(𝐴) ⊂ 𝑈} where 𝐴 ⊂ 𝑋 is compact and 𝑈 ⊂ 𝑌 is
open is called the compact-open topology on 𝐶 𝑋, 𝑌 . The compact-open topology on the
set 𝐶(𝑋, 𝑌) of all continuous maps between two spaces 𝑋 and 𝑌 is defined as follows:
Given a compact subset 𝐾 of 𝑋 and an open subset 𝑈 of 𝑌, let 𝑉(𝐾, 𝑈) denote the set of
all maps 𝑓 in 𝐶(𝑋, 𝑌) such that 𝑓(𝐾) is contained in 𝑈. Then the collection of all
such 𝑉(𝐾, 𝑈) is a subbase for the compact-open topology.

COMPACT OPERATOR: T is a compact operator if whenever (𝑥𝑖 )1 ∞ is a bounded

sequence in 𝑋 then its image (𝑇𝑥𝑖 )1 ∞ has a convergent subsequence in 𝑌.
The set of finite rank operators is denote by 𝐹(𝑋, 𝑌) and the set of compact operators—
by (𝑋, 𝑌) . Both 𝐹(𝑋, 𝑌) and 𝐾(𝑋, 𝑌) are linear subspaces of 𝐵(𝑋, 𝑌). For any two metric
spaces 𝑋 and 𝑌 we have 𝐹(𝑋, 𝑌) ⊂ 𝐾(𝑋, 𝑌).
Let 𝑇𝑚 be a sequence of compact operators convergent to an operator 𝑇 in the norm
topology (𝑖. 𝑒. ||𝑇 − 𝑇𝑚 || → 0) then 𝑇 is compact itself. Equivalently 𝐾(𝑋, 𝑌) is a closed
subspace of 𝐵(𝑋, 𝑌).

COMPACT SET IN A METRIC SPACE: A compact set in a metric space is defined by the
property that any its covering by a family of open sets contains a subcovering by a finite
subfamily. In the finite dimensional vector spaces ℝn or ℂn there is the following
equivalent definition of compactness:

 E is bounded and closed;

 E is compact;
 Any infinite subset of E has a limiting point belonging to E.

COMPACT SPACE: A topological space in which any collection of open sets whose union
is the whole space (open cover of space) has a finite number of open sets from this
collection (finite subcover of the cover) whose union is also the whole space.
A topological space is compact when
(i) Each open cover contains a finite subcover OR
(ii) Every convergent net has a convergent subnet OR
(iii) Every filter on X has a convergent refinement OR
(iv) Every ultrafilter converges to at least one point.
COMPARATIVE STATICS: Comparative static is concerned with the dependencies of
optimal solutions on the parameter.

COMPARATIVE STATICS, MONOTONE: Monotone comparative static is concerned the

optimal solutions varying monotonically with the parameter.

COMPARISON OF FILTERS: If 𝐹 and 𝐹0 are two filters on the same set 𝐸, 𝐹 is said to be
finer than 𝐹0 , or 𝐹0 is coarser than 𝐹, if the fundamental family 𝐹 is finer than 𝐹0 , i.e. if
every 𝐴0 ∇ 𝐹0 contains an 𝐴 ∇ 𝐹. If 𝐹 is finer than 𝐹0 and 𝐹0 is finer than 𝐹, 𝐹0 and 𝐹
are said to be equivalent.
∞
COMPARISON TEST: If 𝑎𝑛 , 𝑏𝑛 ∇ 𝑅 and 0 ≤ 𝑎𝑛 ≤ 𝑏𝑛 then whenever 𝑛=1 𝑏𝑛
∞
converges 𝑛=1 𝑎𝑛 must converge.
COMPATIBLE ATLASES: Two m-dimensional smooth atlases on 𝑀 are said to be
compatible, if every chart from one atlas has smooth transition on its overlap with every
chart from the other atlas (or equivalently, if their union is again an atlas). It can be seen
that compatibility is an equivalence relation. An equivalence class of smooth atlases is
called a smooth structure. All atlases on a given manifold 𝑆 in 𝑅 𝑛 are compatible. The
smooth structure so obtained on 𝑆 is called the standard smooth structure.
COMPATIBLE MATRICES: Two matrices in a particular order so that they can be
multiplied. In the usual convention, the number of columns in the first matrix and the
number of rows in the second matrix must be the same. In more abstract terms,
considering the matrices as transformations, it amounts to maintaining that the number
of dimensions (rank) of the co-domain of the first transformation be the same as the
dimensions (rank) of the domain of the second transformation.

COMPETITIVE GAME: A game is said to be competitive game if it has the following four
properties:

1. There should be finite number of players means 𝑛 ≥ 2.

2. Each player has a finite list of his possible course of action.
3. A play is said to be played when each player choose one of his course of action
and no player knows the choice of action of the other player until he has decided his
own.
4. When each player chooses his activity, then this combination of activities gives a
result according which player gains a payment which may be –ve, +ve or zero.

COMPLEMENT: Let 𝐴 be a subset of some universal set 𝐸. Then the complement of 𝐴 is

the difference set 𝐸\𝐴 (𝑜𝑟 𝐸 – 𝐴). It may be denoted by 𝐴′ (or Ā) when the universal set
is understood or has previously been specified. Complementation is a unary operation
on the set of subsets of a universal set E.
The following properties hold:
 𝐸′ = Ø and Ø′ = 𝐸.
 For all 𝐴, (𝐴′)′ = 𝐴.
 For all 𝐴, 𝐴 ∩ 𝐴′ = Ø and 𝐴 ∪ 𝐴′ = 𝐸

COMPLEMENTARY ANGLES: Two angles that sum to a right angle. In this case, each
angle is called the complement of the other angle.

COMPLEMENTARY EVENTS (PROBABILITY): Events which are both exhaustive and

mutually exclusive. So for A and B to be complementary, 𝑃(𝐴 ∪ 𝐵) = 1 and
𝑃(𝐴 ∩ 𝐵) = 0. It then follows that 𝑃(𝐵) = 1 – 𝑃(𝐴).

COMPLEMENTARY FUNCTION: Along with the particular integral, it forms the general
solution of a linear differential equation. It is essentially an element in the kernel of
the differential operator.

COMPLEMENTARY SUBSPACES: Two subspaces 𝑊1 𝑎𝑛𝑑 𝑊2 of V are called

complementary if 𝑊1 ∩ 𝑊2 = 0 and 𝑊1 + 𝑊2 = 𝑉

Let 𝑊1 , 𝑊2 are subspaces of V. Then 𝑊1 , 𝑊2 are complementary subspaces iff each

vector in v 𝜖 𝑉 can be written uniquely as v= 𝑤1 + 𝑤2 with 𝑤1 𝜖 𝑊1 , 𝑤2 𝜖 𝑊2 .

Examples: 1. Let V=ℝ2 , let 𝑊1 = 𝛼, 0 : 𝛼 𝜖ℝ 𝑎𝑛𝑑 𝑊2 = 0, 𝛼 : 𝛼 𝜖ℝ .Then 𝑊1 , 𝑊2 are

complementary subspaces.
2. V=ℝ3 , 𝑊1 = 𝛼, 0,0 : 𝛼 𝜖ℝ 𝑎𝑛𝑑 𝑊2 = 0, 𝛼, 𝛽 : 𝛼, 𝛽 𝜖ℝ . Here 𝑊1 , 𝑊2 are
complementary subspaces.

3. Let V=ℝ2 , let 𝑊1 = 𝛼, 𝛼 : 𝛼 𝜖ℝ 𝑎𝑛𝑑 𝑊2 = −𝛼, 𝛼 : 𝛼 𝜖ℝ .Then 𝑊1 , 𝑊2 are

complementary subspaces.

COMPLETE ANALYTIC FUNCTION: Suppose 𝑓(𝑧) is analytic in a domain 𝐷. Let us form

all possible analytic continuation of (𝑓, 𝐷) and then all possible analytic continuations of
𝑓1 , 𝐷1 , 𝑓2 , 𝐷2 , 𝑓3 , 𝐷3 , … (𝑓𝑛 , 𝐷𝑛

and so on. At some stage we arrive at a function 𝐹(𝑧) such that for any 𝑣, 𝐹(𝑣) denotes
the value of values obtained for 𝑣 by all possible continuation to 𝑣, that is to say.

𝑓1 𝑧 𝑖𝑓 𝑧 ∇ 𝐷1
𝑓2 𝑧 𝑖𝑓 𝑧 ∇ 𝐷2
𝐹 𝑧 =
…………………
𝑓𝑛 𝑧 𝑖𝑓 𝑧 ∇ 𝐷𝑛

Such a function 𝑓(𝑧) is called complete analytic function. In this process of continuation,
we may arrive at a closed curve beyond which it is not possible to take analytic
continuation. Such a closed curve is called the natural boundary of the complete analytic
function. A point outside the natural boundary is called the singularity of complete
analytic function.

COMPLETE INDUCTION: Also known as strong induction. This method assumes that the
statement is true for all values below a certain finite value in the inductive step of
proving the next statement. Logically strong induction is equivalent to weak induction
and is not "stronger" in this logical sense.

COMPLETE LATTICE: A poset where all subsets have a supremum and an infimum.

COMPLETELY REGULAR SPACE: A space is said to be completely regular (also called a

𝑇31 -space) if it is a 𝑇1 -space and for each point 𝑥 and each closed set 𝐴 with 𝑥 ∇ 𝐴𝑐
2

there is a map 𝑓 ∇ 𝐶(𝑋) such that 𝑓(𝑥) = 𝑎 and 𝑓(𝐴) = {𝑏} where 𝑎 ≠ 𝑏. Thus in a
completely regular space a point and a closed set disjoint from it can be separated by a
continuous real function.
COMPLETE GRAPH: A simple graph in which every vertex is joined to every other. The
complete graph with 𝑛 vertices, denoted by 𝐾𝑛 , is regular of degree 𝑛 – 1 and has edges.

𝐾4 𝐾5 𝐾6

COMPLETE METRIC SPACE: A metric space is said to be a complete metric space in

which every Cauchy sequence is convergent. For example, the real numbers, with the
usual metric.
COMPLETENESS: A system S of numbers is said to be complete if every non-empty
subset S, which is bounded above has a member of S for its supremum. COMPLETE
ORDERED FIELD: An ordered field F is said to be a complete ordered field if every non-
empty subset S of F which is bounded above has an element of F for its supremum. The
set of all real numbers is a complete ordered field.
COMPLETE ORTHONORMAL SEQUENCE: An orthonormal sequence ⌌𝑒𝑛 ⌍1∞ in a Hilbert
space 𝐻 is complete if the identities ⟨ 𝑦, 𝑒𝑘 ⟩ = 0 for all 𝑘 imply 𝑦 = 0. A complete
orthonormal sequence is also called orthonormal basis in H.
COMPLETE SET OF RESIDUES (modulo n): A set of 𝑛 integers, one from each of the 𝑛
residue classes modulo 𝑛. Thus {0, 1, 2, 3, 4} is a complete set of residues modulo 5; so
too are {1, 2, 3, 4, 5} and {−2, – 1, 0, 1, 2}.
COMPLETE SOLUTION OF A DIFFERENTIAL EQUATION: The solution of differential
equation containing a particular integral and the complementary function. This
generates a family of solutions which contains all possible solutions.

COMPLETING THE SQUARE: It is the method for solving a quadratic equation in general
by writing a quadratic expression in the vertex form. Note that the quadratic formula is
derived from the method of completing the square for the general quadratic expression.

COMPLETION: Let 𝑋, 𝜌 be an arbitrary metric space. The complete metric space

𝑋′, 𝜌′ is said to be a completion of 𝑋, 𝜌 if

(i) 𝑋, 𝜌 is isometric to a sub-space 𝑋0 , 𝜌′ of 𝑋 ′ , 𝜌′ , and

 (ii) The closure of 𝑋0 i.e. 𝑋0 , is all of 𝑋 ′ , 𝑖. 𝑒. 𝑋0 = 𝑋 ′ .

An equivalent way of stating condition (ii) is that 𝑋0 is dense of 𝑋 ′ . It means that every
point of 𝑋 ′ is either a point or a limit point of 𝑋0, which in turn implies that, given any 𝑓
of 𝑋 ′ , there exists a sequence of points in 𝑋0 which converges to 𝑓.

COMPLEX ANALYSIS: Complex analysis is the area of mathematics relating to the study
of complex functions.
COMPLEX CONJUGATE: Given a complex number, the complex conjugate is the complex
number whose real part is the same, while the imaginary part (being a real number)
has opposite signs. The significance of complex conjugates stems from the theorem that
says the complex conjugates of all roots of real polynomials are also roots themselves.

COMPLEX FRACTION: A fraction consisting of complex numbers. Considering

the division of complex numbers as complex fractions is a standard way of calculating
the division. (through algebraic methods such as the difference of two squares.)

COMPLEX FUNCTION: Complex function is a function involving complex variables as

either input or output, but usually both. So if 𝑧 = 𝑥 + 𝑦𝑖, 𝑓 𝑧 = 𝑧 2 , 𝑔 𝑧 =
sin 𝑧 , 𝑕 𝑧 = 𝑧 etc are complex functions.
COMPLEX LINE INTEGRAL: Suppose 𝑓(𝑧) is continuous at every point of a closed curve
𝐶 having a finite length i.e. 𝐶 is a rectifiable curve.

Divide 𝐶 into 𝑛 parts by means of points

𝑧0 , 𝑧1 , 𝑧2 , … , 𝑧𝑛 ,

Let 𝑎 = 𝑧0 , 𝑏 = 𝑧𝑛

We choose a point 𝜉𝑘 on each are joining 𝑧𝑘−1 𝑡𝑜 𝑧𝑘 .

Form the sum

𝑆𝑛 = 𝑓 𝜉𝑘 (𝑧𝑟 − 𝑧𝑟 − 1)
𝑟=1

Suppose maximum value of 𝑧𝑟 − 𝑧𝑟 − 1 ⟶ 0 𝑎𝑠 𝑛 ⟶ ∞

Then the sum 𝑆𝑛 tends to a fixed limit which does not depend upon the mode of
subdivision and denote this limit by

𝑏
𝑓 𝑧 𝑑𝑥 𝑜𝑟 ∫𝑐 𝑓 𝑧 𝑑𝑧
𝑧

Which is called the complex line integral or line integral of 𝑓(𝑧) along 𝐶. An evaluation
of integral by such method is also called ab-initio method.

COMPLEX NUMBER: There is no real number 𝑥 satisfying the quadratic equation

𝑥 2 + 1 = 0. The introduction of the imaginary number 𝑖 such that 𝑖 2 = – 1 gives rise
to further numbers of the form 𝑎 + 𝑏𝑖. A number of the form 𝑎 + 𝑏𝑖, where 𝑎 and 𝑏 are
real, is a complex number. Since one may take b = 0, this includes all the real numbers.
The set of all complex numbers is usually denoted by 𝑪. Thus the set C of complex
numbers is closed under addition and multiplication, and the elements of this enlarged
number system satisfy the laws commonly expected of numbers.
It can be verified that addition and multiplication are associative and commutative, that
the distributive law holds, that there is a zero element and an identity element and that
every element has a negative and every non-zero element has an inverse. This shows
that 𝑹 × 𝑹 with this addition and multiplication is a field whose elements, according to
this approach, are called complex numbers. The elements of the form (𝑎, 0) can be seen
to behave exactly like the corresponding real numbers 𝑎.
COMPLEX PLANE: Let points in the plane be given coordinates (𝑥, 𝑦) with respect to a
Cartesian coordinate system. The plane is called the complex plane when the point
(𝑥, 𝑦) is taken to represent the complex number 𝑥 + 𝑦𝑖.
COMPLEX POTENTIAL AND COMPLEX VELOCITY: Since the function ∅(𝑥, 𝑦) and Ψ(𝑥, 𝑦)
constitute the Cauchy. Riemann equation so they provide the necessary and sufficient
condition for the function

𝐹 𝑧 = ∅ 𝑥, 𝑦 + 𝑖Ψ(𝑥, 𝑦)

To be an analytic function of the complex variable 𝑧 = (𝑥 + 𝑖𝑦). Thus the real and
imaginary parts of any analytic function may be regarded as the potential function and
stream function of a flow of an inviscid fluid in two dimensions. The complex function 𝐹,
whose real and imaginary parts are the velocity potential and stream function
respectively, is termed the complex potential of the liquid motion. It ceases to exist at
exist at point occupied by sources, sinks to vortices. Differentiating (1) partially with
regard to 𝑥, we have

dF ∂z ∂∅ ∂Ψ
∙ = +i
dz ∂x ∂x ∂x

dF ∂F
=
dz ∂z

as F is a function of z only .

or
dF ∂∅ ∂∅ ∂Ψ ∂∅
= −i = −𝑢 + 𝑖𝑣; =
dz ∂x ∂y ∂x ∂y

which is called the complex velocity.

COMPLEX VARIABLE: If a simple z takes any one of the values of a set of complex
numbers, then z is called a complex variable.

COMPONENT (CONNECTED SPACES, TOPOLOGY): A connected set in a topological space

is said to be a component if it is not contained in any other connected set of that
topological space other than the topological space itself.
COMPONENT (GRAPH): A graph may have several small pieces and these are called its
components: two vertices are in the same component if and only if there is a path from
one to the other. A more precise definition can be given by defining an equivalence
relation on the set of vertices with u equivalent to v if there is a path from u to v. Then
the components are the corresponding equivalence classes.
COMPONENT (VECTOR): In a Cartesian coordinate system in 3-dimensional space, let
𝒊, 𝒋 and 𝒌 be unit vectors along the three coordinate axes. Given a vector 𝒑, there are
unique real numbers 𝑥, 𝑦 and 𝑧 such that 𝒑 = 𝑥𝒊 + 𝑦𝒋 + 𝑧𝒌. Then 𝑥, 𝑦 and 𝑧 are called
as the components of 𝒑.
COMPOSITE NUMBER: A positive integer is composite if it is neither prime, nor equal to
1; that is, if it can be written as a product 𝑕𝑘, where the integers 𝑕 and 𝑘 are both
greater than 1.
COMPOSITION LAW: If (𝑋, 𝑑) and (𝑌, 𝜌) and (𝑍, 𝜍) are metric spaces and 𝑔 ∶ 𝑋 → 𝑌 , 𝑓 ∶
𝑌 → 𝑍 are continuous, then so is the composition 𝑓 ∘ 𝑔.
COMPOSITION OF FUNCTIONS: Let 𝑓: 𝑆 → 𝑇 and 𝑔: 𝑇 → 𝑈 be mappings. With each
𝑠 ∇ 𝑆 is associated the element 𝑓(𝑠) ∇ 𝑇, and hence the element 𝑔(𝑓(𝑠)) ∇ 𝑈. This
rule gives a mapping from 𝑆 to 𝑈, which is denoted by 𝑔°𝑓 and is the composition of 𝑓
and 𝑔. Thus 𝑔°𝑓: 𝑆 → 𝑈 is defined by (𝑔 °𝑓)(𝑠) = 𝑔(𝑓(𝑠)), and exists if and only if the
domain of 𝑔 equals the range of 𝑓.
COMPOSITION OF LINEAR MAPS: Let 𝑇1 : 𝑈 → 𝑉 𝑎𝑛𝑑 𝑇2 : 𝑈 → 𝑉 be two linear maps. We
define a map 𝑇2 𝑇1 : 𝑈 → 𝑊 by

𝑇2 𝑇1 𝑢 = 𝑇2 𝑇1 𝑢 ∀𝑢 𝜖 𝑈

In particular, we define 𝑇 2 = 𝑇𝑇 and 𝑇 𝑖+1 = 𝑇 𝑖 𝑇 for 𝑖 > 2

It is easy to check that 𝑇1 + 𝑇2 , 𝛼𝑇 and 𝑇2 𝑇1 are themselves all linear maps.

COMPOUND DISTRIBUTION: The compound probability distribution is the result of a

probability distribution whose parameters are distributed along other probability
distributions.

COMPOUND FRACTION: A fraction that consists of another fraction as

its numerator or denominator.

COMPOUND INTEREST: The calculation of interest payments taking into account of

previous interest payments as part of the principal.

COMPOUND STATEMENT: A statement formed from simple statements by the use of

words such as ‘and’, ‘or’, ‘not’, ‘implies’ or their corresponding symbols. The simple
statements involved are the components of the compound statement. For example,
(𝑝 ∧ 𝑞) ∨ (¬𝑟) is a compound statement built up from the components 𝑝, 𝑞 and 𝑟.
COMPRESSION: A name for a transformation where a figure becomes
proportionally smaller.

COMPUTATIONAL FLUID DYNAMICS: Computational fluid dynamics (CFD) is the use of

applied mathematics, physics and computational software to visualize how
a gas or liquid flows -- as well as how the gas or liquid affects objects as it flows past.
Computational fluid dynamics is based on the Navier-Stokes equations. These equations
describe how the velocity, pressure, temperature, and density of a moving fluid are
related.

Computational fluid dynamics has been around since the early 20th century and many
people are familiar with it as a tool for analyzing air flow around cars and aircraft. As
the cooling infrastructure of server rooms has increased in complexity, CFD has also
become a useful tool in the data center for analyzing thermal properties and modeling
air flow. CFD software requires information about the size, content and layout of the
data center. It uses this information to create a 3D mathematical model on a grid that
can be rotated and viewed from different angles. CFD modeling can help an
administrator identify hot spotsand learn where cold air is being wasted or air is
mixing.

Simply by changing variables, the administrator can visualize how cold air will flow
through the data center under a number of different circumstances. This knowledge can
help the administrator optimize the efficiency of an existing cooling infrastructure and
predict the effectiveness of a particular layout of IT equipment. For example, if an
administrator wanted to take one rack of hard drive storage and split the hard drives
over two racks, a CFD program could simulate the change and help the administrator
understand what adjustments would be need to be made to deal with the additional
heat load before any time or money has been spent.

CONCAVE: A geometric figure where it is possible to form a line between two points in
the figure where the line consists of points not from the figure. For a plane figure, it
is equivalent to a shape having an interior angle of greater than 180 degrees.

CONCAVE FUNCTION: A function whose graph is such that, for any two points of the
graph, the function for arguments between the two points are higher than the straight
line joining the two points. For a differentiable function, it is equivalent to a function
with a monotonically decreasingly gradient. Note that a function can be described
as concave for a certain interval only.

CONCAVE POLYGON: A polygon, as a plane figure, which is concave.

CONCAVITY: At a point of a graph 𝑦 = 𝑓(𝑥), it may be possible to specify the concavity
by describing the curve as either concave up or concave down at that point, as follows, If
the second derivative 𝑓″(𝑥) exists and is positive throughout some neighbourhood of a
point 𝑎, then 𝑓′(𝑥) is strictly increasing in that neighbourhood, and the curve is said to
be concave up at 𝑎. At that point, the graph 𝑦 = 𝑓(𝑥) and its tangent look like one of the
cases shown in the first figure.

If 𝑓″(𝑎) > 0 and 𝑓″ is continuous at a, it follows that 𝑦 = 𝑓(𝑥) is concave up at 𝑎.

Consequently, if 𝑓′(𝑎) = 0 and 𝑓″(𝑎) > 0, the function 𝑓 has a local minimum at 𝑎.
Similarly, if 𝑓″(𝑥) exists and is negative throughout some neighbourhood of 𝑎, or if
𝑓″(𝑎) < 0 and 𝑓″ is continuous at 𝑎, then the graph 𝑦 = 𝑓(𝑥) is concave down at 𝑎 and
looks like one of the cases shown in the second figure. If 𝑓′(𝑎) = 0 and 𝑓″(𝑎) < 0, the
function 𝑓 has a local maximum at 𝑎.
CONCAVE SET: A set of three or more points is concave if it is possible to draw a line
segment that connects two points that are in the set, but includes also some points that
are not in the set.

CONCENTRIC CIRCLES: A collection of circles is said to be concentric if they have the

same center. The region formed between two concentric circles is therefore an annulus.
CONCURRENT: The property of sharing a common point.

CONCYCLIC POINTS: A number of points are said to be concyclic points if there is a circle
that passes though all of them.
CONDITIONAL DISTRIBUTION: The distribution of a random variable given that another
random variable is known to be of a certain value.

CONDITIONAL EQUATION: An equation which is only true under certain contexts.

Normally known simply as equations, as opposed to identities, which is true regardless
of contexts.

CONDITIONAL INEQUALITY: An inequality that is only true under certain conditions.

CONDITIONAL STATEMENT: A statement of the type that something will be true

provided that something else is true. For example ‘if 𝑛 is not divisible by 2, 𝑛 is odd is a
conditional statement.
CONDITIONAL PROBABILITY: For two events 𝐴 and 𝐵, the probability that 𝐴 occurs,
given that 𝐵 has occurred, is denoted by 𝑃(𝐴|𝐵), read as ‘the probability of 𝐴 given 𝐵’.
This is called a conditional probability. Provided that 𝑃(𝐵) is not zero, 𝑃(𝐴|𝐵) =
𝑃(𝐴 ∩ 𝐵)/𝑃(𝐵). This result is often useful in the following form: 𝑃(𝐴 ∩ 𝐵) =
𝑃(𝐵) 𝑃(𝐴|𝐵). If 𝐴 and 𝐵 are independent events, 𝑃(𝐴|𝐵) = 𝑃(𝐴), and this gives the
product law for independent events: 𝑃(𝐴 ∩ 𝐵) = 𝑃(𝐴) 𝑃(𝐵).
CONDITIONALY CONVERGENT SERIES: A series 𝑢𝑛 is said to be absolutely convergent
if 𝑢𝑛 is convergent but the series 𝑢𝑛 is not convergent.
CONDITION FOR CONSISTENCY: The system of equations 𝐴𝑋 = 𝐵 is consistent i.e.,
possesses a solution, if and only if the coefficient matrix 𝐴 and the augmented matrix
[𝐴, 𝐵] are of the same rank.

CONDITION FOR FOUR POINTS TO BE COPLANAR: The necessary and sufficient

condition for any four points in three dimensional space to be coplanar is that there
exists a linear relation between their position vectors such that the algebraic sum of the
scalar coefficients in it is zero, provided the scalars are nor all zero.

CONDITION FOR THREE POINTS TO BE COLLINEAR: The necessary and sufficient

condition for three points in three dimensional space to be collinear is that there exists
a linear relation connecting their position vectors such that the algebraic sum of the
scalar coefficients in it is zero, provided the scalars are not all zero.
CONDITION, NECESSARY AND SUFFICIENT (DISCRETE MATHEMATICS): The
implication 𝑞 ⇒ 𝑝 can be read as ‘if 𝑞 then 𝑝’. When this is true, it may be said that 𝑞 is
a sufficient condition for 𝑝; that is, the truth of the ‘condition’ 𝑞 is sufficient to ensure
the truth of 𝑝. This means that 𝑝 is true if 𝑞 is true. On the other hand, when the
implication 𝑝 ⇒ 𝑞 holds, then 𝑞 is a necessary condition for 𝑝; that is, the truth of the
‘condition’ 𝑞 is a necessary consequence of the truth of 𝑝. This means that 𝑝 is true only
if 𝑞 is true. When the implication between 𝑝 and 𝑞 holds both ways, 𝑝 is true if and only
if q is true, which may be written 𝑝 ⇔ 𝑞. Then 𝑞 is a necessary and sufficient condition
for 𝑝.
CONFIDENCE INTERVAL (CONFIDENCE LEVEL): An interval, calculated from a sample,
which contains the value of a certain population parameter with a specified probability.
The end-points of the interval are the confidence limits. The specified probability is
called the confidence level. An arbitrary but commonly used confidence level is 95%,
which means that there is a one-in-twenty chance that the interval does not contain the
true value of the parameter. For example, if 𝑥 is the mean of a sample of 𝑛 observations
taken from a population with a normal distribution with a known standard deviation 𝜍,
1.96 𝜍 1.96 𝜍
then 𝑥 − , 𝑥+ is a 95% confidence interval for the population mean 𝜇.
𝑛 𝑛

CONFIDENCE LIMITS: Two values between which the true value of a population
parameter is thought to lie with some given probability. This probability represents the
proportion of occasions when such limits calculated from repeated samples actually
include the true value of the parameter. An essential feature of the interval is that the
distance between the limits depends on the size of the sample, being smaller for a larger
sample. CONFOCAL CONICS: Two central conics are said to be confocal conics if they
have the same foci. An ellipse and a hyperbola that are confocal intersect at right angles.
CONFOCAL QUADRICS: A family of central quadrics represented by the following
equations is called a family of confocal quadrics.

𝑥2 𝑦2 𝑧2
+ 𝑏+𝑘 + 𝑐+𝑘 = 1, 𝑎 > 𝑏 > 𝑐 > 0,
𝑎+𝑘

Where 𝑘 is a parameter. For a quadric belonging to this family, any point on the ellipse
𝑥2 𝑎 − 𝑐 + 𝑦2 𝑏 − 𝑐 = 1, 𝑧 = 0 or the hyperbola
𝑥2 𝑎 − 𝑏 − 𝑧2 𝑏 − 𝑐 = 1, 𝑦 = 0 is called a focus. The ellipse and hyperbola are
called focal conics of the quadric.
CONFORMABLE MATRICES: Two matrices A and B are said to be conformable matrices
(for multiplication) if the number of columns of A equals the number of rows of B. Then
A has order 𝑚 × 𝑛 and B has order 𝑛 × 𝑝, for some 𝑚, 𝑛 and 𝑝, and the product AB, of
order 𝑚 × 𝑝, is defined.

CONFORMALITY: Now we make some remarks about “conformality.” A function is

conformal at a point 𝑃 ∇ 𝐶 if the function “preserves angles” at 𝑃 and “stretches
equally in all directions” at 𝑃. Both of these statements must be interpreted
infinitesimally. Holomorphic functions enjoy both properties.
CONFORMALLY FLAT: An 𝑀 is conformally flat if it is locally conformally equivalent to a
Euclidean space, for example standard sphere is conformally flat.

CONFORMAL MAPPING (DIFFERENTIAL GEOMETRY): A surface S is said to be

conformally mapped on to a surface S* if there is a differentiable homomorphism (i.e.
and analytic one-one map) of S on to S* such that the angle between any two curves at
any arbitrary point P on S is equal to the angle between the corresponding directed
curves on S*. Thus in short, a map said to be conformal if it preserves angles.

CONFORMAL TRANSFORMATION (COMPLEX ANALYSIS): Suppose the transformation

𝑢 = 𝑢 𝑥, 𝑦 , 𝑣 = 𝑣(𝑥, 𝑦) maps the two curves 𝑐1 , 𝑐2 [intersecting at the point p 𝑤0 of 𝑤
plane .

If the angle between 𝑐1 and 𝑐2 at 𝑧0 is equal to the angle between 𝑐1 and 𝑐2 at 𝑤0 then the
transformation is called isogonal. If the sense of rotation as well as the magnitude of the
angle is preserved the transformation is called conformal.

CONGRUENCE (modulo n): For each positive integer 𝑛, the relation of congruence
between integers is defined as follows: 𝑎 is congruent to 𝑏 modulo 𝑛 if 𝑎 – 𝑏 is a
multiple of 𝑛. This is written as 𝑎 ≡ 𝑏 (𝑚𝑜𝑑 𝑛). The integer 𝑛 is the modulus of the
congruence. Then 𝑎 ≡ 𝑏 (𝑚𝑜𝑑 𝑛) if and only if 𝑎 and 𝑏 have the same remainder upon
division by 𝑛. For example, 33 is congruent to 8 modulo 5.
The following properties hold, if 𝑎 ≡ 𝑏 (𝑚𝑜𝑑 𝑛) and 𝑐 ≡ 𝑑 (𝑚𝑜𝑑 𝑛):
 𝑎 + 𝑐 ≡ 𝑏 + 𝑑 (𝑚𝑜𝑑 𝑛),
 𝑎 – 𝑐 ≡ 𝑏 – 𝑑 (𝑚𝑜𝑑 𝑛),
 𝑎𝑐 ≡ 𝑏𝑑 (𝑚𝑜𝑑 𝑛).
It can be shown that congruence modulo n is an equivalence relation and so defines a
partition of the set of integers, where two integers are in the same class if and only if
they are congruent modulo n. These classes are the residue (or congruence) classes
modulo n.
CONGRUENCE EQUATION: congruence equations are the equations involving
congruence modulo n. The following are examples of congruence equations:
(i) 𝑥 + 5 ≡ 3 (𝑚𝑜𝑑 7); this has the solution 𝑥 ≡ 5 (𝑚𝑜𝑑 7).
(ii) 2𝑥 ≡ 5 (𝑚𝑜𝑑 4); this has no solutions.
In seeking solutions to a congruence equation, it is necessary only to consider a
complete set of residues and find solutions in this set.
The linear congruence equation 𝑎𝑥 ≡ 𝑏 (𝑚𝑜𝑑 𝑛) has a solution if and only if (𝑎, 𝑛)
divides 𝑏, where (𝑎, 𝑛) is the greatest common divisor of 𝑎 and 𝑛.
CONGRUENCE OF MATRICES: A square matrix B of order 𝑛 over a field 𝐹 is said to be
congruent to another square matrix 𝐴 of order 𝑛 over 𝐹, if there exists a non-singular
matrix 𝑃 over 𝐹 such that 𝐵 = 𝑃 −1 𝐴𝑃.

 The relation of ‘congruence of matrices’ is an equivalence relation in the set of all

n × n matrices over a field 𝐹.
 Every matrix congruent to a symmetric matrix is a symmetric matrix.
 Each congruent transformation of a matrix consists of a pair of elementary
transformations, one row and the other columns, such that of the corresponding
elementary matrices each is the transpose of the other.
 Every matrix B obtained from any given matrix 𝐴 by subjecting 𝐴 to a finite chain
of congruent operations is congruent to 𝐵.

CONGRUENCE OF QUADRATIC FORMS OR EQUIVALENCE OF QUADRATIC FORMS: Two

quadratic forms X T AX and Y T BY over a field F are said to be congruent or equivalent
over F if their respective matrices A and B are congruent over F. Thus X T AX is
equivalent to Y T BY if there exists a non-singular matrix P over F such that 𝑃 −1 𝐴𝑃 = 𝐵.
Since congruence of matrices is an equivalence relation therefore equivalence of
quadratic forms is also an equivalence relation.
CONGRUENT: One figure that would coincide with another with
a combination of translation, rotation and reflection. Essentially, two shapes are
congruent if they can be considered "the same" except for its location and orientation.

For triangles, there are certain conditions that make it easier to decide if two such
figures are the same. (Without having to consider the transformations)

I - SSS - The length of all three sides of a triangle are the same as the lengths of the
corresponding sides of another.

II - SAS - The lengths of any two pairs of sides correspond, and the angles between those
sides are also the same as each other.

III - ASA/AAS - Any pair of angles on one shape correspond with the other shape, while
any side of either shape is the same as the corresponding side of the other.

IV - RHS (also known as LH) - Both triangles are right-angled, their hypotenuses are of
the same length and they share another side of the same length.

CONICAL PENDULUM: A weight attached through a string to a fixed point so that

the trajectory of the weight is a (horizontal) circle with the string being taut (and
having constant length) at all times. The circle drawn out by the weight together with
the positions of the string forms a cone which explains the name.

CONICOID: A surface generated by the rotation about an axis of one of the conics -
ellipsoids, paraboloids, hypocycloids and a sphere.

CONIC SECTIONS: It is a curve produced when a plane intersects a right-circular cone.

The four curves—circles, ellipses, parabolas, and hyperbolas are called conic sections
because they can be formed by the intersection of a plane with a right circular cone. If
the plane is perpendicular to the axis of the cone, the intersection will be a circle. If the
plane is slightly tilted, the result will be an ellipse. If the plane is parallel to one element
of the cone, the result will be a parabola. If the plane intersects both parts of the cone,
the result will be a hyperbola.
 b

CONJECTURE: An assertion that is not yet proven. In this sense, it is the same as
an hypothesis.

CONJUGACY: Two elements 𝑕 and 𝑘 of a group 𝐺 are said to be conjugate if 𝑘 = 𝑔𝑕𝑔−1

for some 𝑔 ∇ 𝐺.
CONJUGACY CLASS: The set of all elements of a group that are conjugate to an element 𝑎
is called cojugacy class of 𝑎.
CONJUGATE (COMPLEX NUMBER): For any complex number 𝑧, where 𝑧 = 𝑥 + 𝑦𝑖, its
conjugate read as ‘z bar’ is equal to 𝑧 = 𝑥 – 𝑦𝑖. In the complex plane, the points
representing a complex number and its conjugate are mirror images with respect to the
real axis. The following properties hold:

It is an important fact that if the complex number 𝛼 is a root of a polynomial equation

𝑧 𝑛 + 𝑎1 𝑧 𝑛–1 + ··· + 𝑎𝑛–1 𝑧 + 𝑎𝑛 = 0, where 𝑎1 , − − − 𝑎𝑛–1 , 𝑎𝑛 are real, then ᾱ is also a

root of this equation.
CONJUGATE AXIS: It is the axis of symmetry perpendicular to the transverse axis of a
hyperbola.
CONJUGATE ELEMENTS: Two elements 𝑥 and 𝑦 in a group 𝐺 are said to be conjugate to

each other if there is an element 𝑎 in 𝐺 for which 𝑦 = 𝑎–1 𝑥𝑎.

CONJUGATE FUNCTIONS: If 𝑓 𝑧 = 𝑢 + 𝑖𝑣 is analytic and if 𝑢 and 𝑣 satisfy Laplace’s
equation ∆2 = V = 0, then 𝑢 and 𝑣 are called conjugate harmonic function or conjugate
functions simply.

CONJUGATE OF A MATRIX: The matrix obtained from any given matrix A on replacing its
elements by the corresponding conjugate complex numbers is called the conjugate of A
and is denoted by A.

Thus if A = aij where aij denoted the conjugate complex of aij .

m×n

If A be a matrix over the field of real numbers, then obviously A coincides with A.

2 + 3i 4 − 7i 8 2 − 3i 4 + 7i 8
A= , then A = .
−i 6 9+i i 6 9−i

If A and B be the conjugates of A and B respectively, then

(i) A = A; (ii) A + B = A + B;
(i) kA = k A; 𝑘 being any complex number;
(ii) AB = A B, A and B being conformable to multiplication.

CONJUGATE POINTS (DIFFERENTIAL GEOMETRY): Two points 𝑝 and 𝑞 on a geodesic 𝛾

are called conjugate if there is a Jacobi field on 𝛾 which has a zero at 𝑝 and 𝑞.

CONJUGATE SETS: Subsets 𝑋 and 𝑌 in a group 𝐺 are said to be conjugate if there is an

element 𝑎 in 𝐺 for which 𝑌 = 𝑎–1 𝑋𝑎.

CONJUGATE SUBGROUPS: Two subgroups 𝐻 and 𝐾 of a group 𝐺 are said to be conjugate
if 𝐾 = 𝑔𝐻𝑔−1 for some 𝑔 ∇ 𝐺.
CONJUNCTION: If 𝑝 and 𝑞 are statements, then the statement ‘𝑝 𝑎𝑛𝑑 𝑞’, denoted by
𝑝 ∧ 𝑞, is the conjunction of 𝑝 and 𝑞. For example, if 𝑝 is ‘It is raining’ and 𝑞 is ‘It is
March’, then 𝑝 ∧ 𝑞 is ‘It is raining and it is March’. The conjunction of 𝑝 and 𝑞 is true
only when 𝑝 and 𝑞 are both true, and so the truth table is as follows:
p Q 𝑝 ∧𝑞
T T T
T F F
F T F
F F F

CONNECTED (RELATION): A binary relation is connected if for all pairs of elements

𝑥, 𝑦, 𝑥 ≠ 𝑦, either 𝑥 ~ 𝑦 or 𝑦 ~ 𝑥. So, for example, in the set of real numbers the relation
‘is greater than’ is connected.
CONNECTED (SET): A set A is said to be a connected set if ∄ non-empty subsets E and F
of A such that 𝐸 ∩ 𝐹 = 𝜑 and 𝐸 ∩ 𝐹 = 𝜑.
CONNECTED COMPONENT: Let 𝛺 be an open set in 𝐶 and 𝑧 ∇ 𝛺. The connected
component (or simply the component) of 𝛺 containing 𝑧 is the set 𝐶𝑧 of all points w in 𝛺
that can be joined to 𝑧 by a curve entirely contained in 𝛺. In other words, a connected
component of a space is a maximal nonempty connected subspace. Each connected
component is closed, and the set of connected components of a space is a partition of
that space.

CONNECTED GRAPH: A connected graph is a graph in which there always exists a path
from any one vertex to any other vertex. So a graph is connected if it is ‘all in one piece’;
that is, if it has precisely one component.
CONNECTED METRIC SPACE: A metric space (𝑋, 𝑑) is said to be connected if there do
not exist two nonempty and disjoint open sets 𝐴 and 𝐵 in 𝑋 such that 𝐴 ∪ 𝐵 = 𝑋. A
metric space (𝑋, 𝑑) is connected if, and only if, the only subsets of 𝑋 that are both open
and closed in 𝑋 are ∅ and 𝑋. Let 𝑋 and 𝑌 be metric spaces, and let 𝑓 ∶ 𝑋 → 𝑌 be a
continuous function. If 𝑋 is connected, then 𝑓(𝑋) is a connected subset of 𝑌.

CONNECTED SET: A set that can't be split into a union of two sets each of which is
both open and closed.
CONNECTED SUM: Let S and T be two surfaces. From each surface, delete an open disk,
then glue the two boundary circles. The resulting surface is called the connected sum of
the two surfaces, denoted by S#T. It is known that the connected sum does not depend
on the choices of the disks.
CONSECUTIVE ANGLES: A set of angles where every member can be considered to be
"next" to another angle within the set. This idea of "next" or adjacency can be sharing
a side or sharing a side as well as a point depending on the context.

CONSECUTIVE SIDES: A set of sides (edges) where every memeber can be considered to
be "next" to another side within the set. This idea of "next" or adjacency can be sharing
a vertex or sharing an angle depending on the context.

CONNECTIVITY: A function 𝑓 ∶ 𝑋 → 𝑌 is connectivity, 𝑓 ∇ 𝐶𝑜𝑛𝑛(𝑋, 𝑌 ), if the graph of

the restriction 𝑓/ 𝑍 of 𝑓 to 𝑍 is connected in 𝑋 × 𝑌 for every connected subset 𝑍 of 𝑋. It
is easy to see that 𝑓 ∶ 𝑅 → 𝑅 is connectivity if and only if its graph is a connected
subset of 𝑅 2 . However, there are functions 𝐹 ∶ 𝑅 2 → 𝑅 with a connected graph which
are not connectivity functions. For example, this is the case if 𝐹 (𝑥, 𝑦) = 𝑠𝑖𝑛(1/𝑥) for
𝑥 = 0, and 𝐹 (0, 𝑦) = 𝑕(𝑦), where 𝑕 ∶ 𝑅 → [−1, 1] is any function with a disconnected
graph.

CONSEQUENT: In the conventional way of expressing hypothetical propositions, "If A

then B", the consequent is the second part of the sentence. It is the part of the sentence
whose truth value is dependent on the other part. (Within the context of the statement
itself)

CONSERVATION LAWS: Any theorem or assertion that states that certain measurable
quantities remains unchanged, and the condition under which this happens.

CONSERVATION OF ENERGY: The law that states that the amount of energy in a closed
system remains unchanged. The exactness of this law, in our current understanding,
depends on whether we use mass-energy as the measure or simply energy in the
classical sense.

CONSERVATION OF MOMENTUM: The law that states that the total momentum of a
closed system of objects remains unchanged.

CONSISTENCY: An axiomatic theory is consistent if it's impossible (in the confines of the
theory) to prove simultaneously a statement and its negation. The Godel's Theorem
states that any (sufficiently powerful) consistent axiomatic theory is incomplete.
CONSISTENT SET OF EQUATIONS: A set of equations is consistent if there is a
solution.
CONSTANT FUNCTION: In real analysis, a constant function is a real function 𝑓 such that
𝑓(𝑥) = 𝑎 for all 𝑥 ∇ 𝑹, where 𝑎, the value of 𝑓, is a fixed real number.
CONSTANT OF INTEGRATION: If 𝜑 is a particular antiderivative of a continuous function
𝑓, then any antiderivative of 𝑓 differs from 𝜑 by a constant. It is common practice,
therefore, to write
∫ 𝑓(𝑥)𝑑𝑥 = 𝜑(𝑥) + 𝑐,
where 𝑐, an arbitrary constant, is the constant of integration.
CONSTANT VALUE THEOREM: If 𝑓 ′ (𝑡) = 0 for all 𝑡 ∇ (𝑎, 𝑏) then 𝑓 is constant on
[𝑎, 𝑏]. (That is, 𝑓(𝑦) = 𝑓(𝑥) whenever 𝑏 ≥ 𝑦 > 𝑥 ≥ 𝑎.)

CONSTRAINTS: The constraints may be defined as that set of linear equations /

inequations under the action of which the objective function is optimized.

CONSTRUCTION OF MEASURES, PRODUCTS: Consider the semiring S of intervals [a,b).

There is a simple description of all measures on it. For a measure µ, define 𝐹𝜇 𝑡 =
µ 0, t if t > 0
0 if t = 0
−µ t, 0 if t < 0

𝐹𝜇 is monotonic and any monotonic function 𝐹 defines a measure µ on 𝑆 by the by

µ([𝑎, 𝑏)) = 𝐹(𝑏) − 𝐹(𝑎). The correspondence is one-to-one with the additional
assumption 𝐹(0) = 0. The above measure µ is 𝜍 −additive on 𝑆 if and only if 𝐹 is
continuous from the left: 𝐹(𝑡 − 0) = 𝐹(𝑡) for all 𝑡 ∇ ℝ.

CONTAINED IN, CONTAINS: It is tempting to say that ‘𝑥 is contained in 𝑆’ when 𝑥 ∇ 𝑆,

and also to say that ‘𝐴 is contained in 𝐵’ if 𝐴 ⊆ 𝐵. To distinguish between these two
different notions, it is better to say that ‘𝑥 belongs to 𝑆’ and to say that ‘𝐴 is included in
𝐵’ or ‘𝐴 is a subset of 𝐵’. However, some authors consistently say ‘is contained in’ for ⊆.
Given the same examples, it is similarly tempting to say that ‘𝑆 contains 𝑥’ and also that
‘𝐵 contains 𝐴.’ It is again desirable to distinguish between the two by saying that
‘𝐵 includes 𝐴’ in the second case, though some authors consistently say ‘contains’ in this
situation. The first case is best avoided or else clarified by saying that ‘𝑆 contains the
element 𝑥’ or ‘𝑆 contains 𝑥 as an element.’
CONTINGENCY TABLE: A table in which individuals are categorized according to two or
more characteristics or variables. Each cell of the table contains the number of
individuals with a particular combination of characteristics or values.
CONTINGENT (LOGIC): A statement or proposition is contingent if it is neither always
true nor always false. For example ‘𝑥 divided by 2 is an integer’ is true when
𝑥 = 2, 4, 6, … but is not true when 𝑥 = 3, 4.5, 7, …
CONTINUED FRACTION: An expression of the form 𝑞1 + 1/𝑏2 , where 𝑏2 = 𝑞2 +
1/𝑏3 , 𝑏3 = 𝑞3 + 1/𝑏4 , and so on, where 𝑞1 , 𝑞2 , 𝑞3 , … are integers, usually positive. If
the continued fraction terminates, it gives a rational number. When the continued
fraction continues indefinitely, it represents an irrational number.

CONTINUITY CORRECTION: The adjustment of the argument of a discrete distribution to

form a closer approximation to a continuous argument. (Usually an adjustment of 0.5
for discrete distributions that admits integers)

CONTINUOUS DISTRIBUTION: A distribution where the domain of the

cumulative distribution function is continuous.

CONTINUOUS FUNCTION: The real function f of one variable is said to be continuous at 𝑎

if 𝑓(𝑥) → 𝑓(𝑎) 𝑎𝑠 𝑥 → 𝑎. The rough idea is that, close to 𝑎, the function has values
close to 𝑓(𝑎). It means that the function does not suddenly jump at 𝑥 = 𝑎 or take
widely differing values arbitrarily close to 𝑎.
A function 𝑓 is continuous in an open interval if it is continuous at each point of the
interval; and 𝑓 is continuous on the closed interval [𝑎, 𝑏], where 𝑎 < 𝑏, if it is
continuous in the open interval (𝑎, 𝑏) and if lim𝑥→𝑎+ 𝑓(𝑥) = 𝑓(𝑎) and lim𝑥→𝑏− 𝑓(𝑥) =
𝑓(𝑏). The following properties hold:
 The sum of two continuous functions is a continuous function.
 The product of two continuous functions is a continuous function.
 The quotient of two continuous functions is continuous at any point or in any
interval where the denominator is not zero.
 Suppose that 𝑓 is continuous at 𝑎, that 𝑓(𝑎) = 𝑏 and that 𝑔 is continuous at 𝑏.
Then 𝑕, defined by 𝑕(𝑥) = (𝑔 𝑓)(𝑥) = 𝑔(𝑓(𝑥)), is continuous at 𝑎.
CONTINUOUS HYPERGEOMETRIC FUNCTIONS: According to Gauss, the function
𝐹 a′ , b′ , c ′ , z′ is continuous to 𝐹 a; b; c; z if it is obtained by increasing or decreasing
one and only one of the parameters a, b, c by unity.

Thus there are six hypergeometric functions continuous to 𝐹 a; b; c; z and are denoted
as

𝐹 a + 1; b; c; z = Fa+, 𝐹 a − 1; b; c; z = Fa−

𝐹 a; b + 1; c; z = Fb+, 𝐹 a; b − 1; c; z = Fb−

𝐹 a; b; c + 1; z = Fc+, 𝐹 a; b; c − 1; z = Fc−

CONTINUOUS MAP (TOPOLOGY): A map 𝑓 ∶ 𝑋 → 𝑌 is continuous if inverse image of an

open set is always open. Note that
(1) A composition of continuous functions is continuous,
(2) 𝑓 continuous ⇒ 𝑓(compact subset) is compact,
(3) 𝑓 continuous ⇒ 𝑓(connected subset) is connected.
(4) Continuous bijection from a compact space to a Hausdorff space ⇒
homeomorphism.
(5) Continuous surjection from a compact space to a Hausdorff space ⇒ quotient map.
CONTINUOUSLY DIFFERENTIABLE FUNCTION: A function is continuously differentiable
if its derivative is itself a continuous function. The function 𝑓(𝑥) = 0 for 𝑥 ≤ 0 and
𝑓(𝑥) = 𝑥 for 𝑥 ≥ 0 is continuous, but not continuously differentiable since 𝑓′(𝑥) = 0
for 𝑥 < 0 and 𝑓′(𝑥) = 1 for 𝑥 > 0.
CONTINUUM: The set of real numbers or any interval (𝑎, 𝑏), which can be open or closed
at either end is a continuum.
CONTINUUM HYPOTHESIS: The conjecture made by Georg Cantor that there is no set
with a cardinal number between aleph-null which is the cardinal number of the set of
natural numbers and the cardinal number of the set of real numbers, i.e. the continuum.
CONTINUOUS RANDOM VARIABLE: A random variable that can take on any real-number
value within a certain finite or infinite interval.
CONTOUR: By counter, we mean a continuous chain of a finite number of regular areas.
If the contour is closed and does not interest itself, then it is called closed contour.

For example, Boundaries of triangles and rectangles are closed contours.

CONTOUR INTEGRAL: An integral ∫ 𝑓(𝑧) 𝑑𝑧 of a function f in the complex plane over a
curve C, usually a closed curve, in the plane.
CONTOUR LINE: A line joining points of a constant value. If 𝑧 = 𝑓(𝑥, 𝑦) is a function
which defines a surface and the line 𝑦 = 𝑔(𝑥) has the property that 𝑓(𝑥, 𝑔(𝑥)) is
constant then the line 𝑦 = 𝑔(𝑥) is a contour line.
CONTRACTION MAPPING: A mapping 𝑓 on a metric space 𝑋 is said to be a contraction
mapping, which decreases the distance, 𝑑, between any two points in 𝑋, i.e. for some
𝑎 < 1, i.e. 𝑑 𝑓(𝑥), 𝑓(𝑦) ≤ 𝑎 × 𝑑 𝑥, 𝑦 for all 𝑥, 𝑦 ∇ 𝑋.
CONTRACTION MAPPING PRINCIPLE: Suppose that 𝑆 is a closed subset of a Banach
space, 𝑌, and that 𝑇 ∶ 𝑆 → 𝑆 is a mapping on 𝑆 such that
||𝑇𝑢 − 𝑇𝑣||𝑌 ≤ 𝛼||𝑢 − 𝑣||𝑌 ; 𝑢, 𝑣 ∇ 𝑆
for some constant 𝛼 < 1. Then T has a unique fixed point 𝑢 ∇ 𝑆 that satisfies
𝑇𝑢 = 𝑢.
CONTRADICTION: The simultaneous assertion of both the truth of a proposition and its
denial. Since both cannot be true there must necessarily be a flaw in either the
reasoning leading to the simultaneous assertion or in the assumptions on which the
deductive reasoning is based. It is this latter situation which provides the basis for proof
by contradiction.
CONTRAPOSITION: The logical principle, upon which proof by contradiction is based.
Let 𝑝 and 𝑞 be statements. If 𝑝 implies not 𝑞, then 𝑞 being true implies that 𝑝 cannot be.
For example, since all rhombuses are parallelograms, a shape which is not a
parallelogram cannot be a rhombus. Here 𝑝 is the statement ‘is a rhombus’ and 𝑞 is the
statement ‘is not a parallelogram’.
CONTRAPOSITIVE: The contrapositive of an implication 𝑝 ⇒ 𝑞 is the implication
¬𝑞 ⇒ ¬𝑝. An implication and its contrapositive are logically equivalent, so that one is
true if and only if the other is. So, in giving a proof of a mathematical result, it may on
occasion be more convenient to establish the contrapositive rather than the original
form of the theorem. For example, the theorem that if 𝑛2 is odd then 𝑛 is odd could be
proved by showing instead that if 𝑛 is even then 𝑛2 is even.
CONTRAVARIANT AND CO-VARIANT VECTORS: Let 𝐴𝑖 , 𝑖 = 1,2 … . 𝑛, be 𝑛 functions of co-
ordinate 𝑥1 , 𝑥 2 , … 𝑥 𝑛 . If the quantities 𝐴𝑖 are transformed to another co-ordinate system
 𝜕𝑥 𝑎
𝑥′1 , 𝑥′2 , 𝑥′3 , … 𝑥′𝑛 according to the rule 𝐴𝑖 = 𝐴𝑎 , then the function Ai are called
𝜕𝑥 𝑖

components of contravariant vector.

Let 𝐴𝑖 , 𝑖 = 1,2 … . 𝑛, be 𝑛 functions of co-ordinates 𝑥1 , 𝑥 2 , … … . , 𝑥 𝑛 . If the quantities 𝐴𝑖

are transformed to another co-ordinate system 𝑥′1 , 𝑥′2 , 𝑥′3 , … 𝑥′𝑛 according to the rule
𝜕𝑥 𝑖
𝐴𝑖 = 𝐴𝑎 , then the functions 𝐴𝑖 are called components of a covariant vector. The
𝜕𝑥 𝑎

contravariant (or covariant) vector is also called a contravariant (or Covariant) tensor
of rank one.

CONVENIENCE SAMPLE: Quick, easy, and readily available people or objects selected to
conduct a survey. It is convenient for the person compiling the statistics. If you conduct
a survey on people´s opinions about a specific movie by interviewing people who just
walked out of that movie, you are taking a convenience sample.
CONVENIENCE SAMPLING: In convenience sampling, a sample is chosen by using the
most conveniently available group. Data collected from such a sample are unlikely to
contain much worthwhile information about a larger population.
CONVERGE (NET): A net 𝑓𝑛 is said to said to converge to a point 𝑥 of the directed set 𝐷,
CONVERGE (SEQUENCE): A real sequence 𝑥𝑛 is said to be convergent to a real number 𝑎
if for every ∇> 0 there exists a 𝑚∇ ∇ 𝑁 such that 𝑥𝑛 − 𝑎 <∇ whenever 𝑛 ≥ 𝑚∇ .
CONVERGE (SEQUENCE IN A TOPOLOGICAL SPACE): A sequence 𝑥𝑛 in a topological
space is said to be convergent to a number 𝑎 if for every for every nbd 𝑁𝑎 of 𝑎 there
exists a 𝑚𝑎 ∇ 𝑁 such that 𝑥𝑛 ∇ 𝑁𝑎 whenever 𝑛 ≥ 𝑚𝑎 .
CONVERGE (SERIES): The infinite series 𝑎1 + 𝑎2 + 𝑎3 + … is said to converge to a
limit 𝐿 if for every 𝜀 > 0 there exists an 𝑁 such that
𝑛
𝑖=1 𝑎𝑛 − 𝐿 <∇ for all 𝑛 > 𝑁.
CONVERGENCE AND DIVERGENCE OF INFINITE SERIES: Let 𝑎𝑛 𝑛 = 1,2,3, … be a
sequence of real or complex numbers. Then the formal infinite sum 𝑎1 + 𝑎2 + ⋯ is
∞
called an infinite series (or series) and is denoted by 𝑛=1 𝑎𝑛 𝑜𝑟 𝑎𝑛 . The number 𝑎𝑛 is
the 𝑛𝑡𝑕 term of the series 𝑎𝑛 , 𝑎𝑛𝑑 𝑠𝑛 = 𝑎1 + 𝑎2 + ⋯ + 𝑎𝑛 is the 𝑛𝑡𝑕 partial sum of
𝑎𝑛 . Also for a finite sequence 𝑎1 , 𝑎2 , … , 𝑎𝑛 , the sum 𝑎1 + 𝑎2 + ⋯ + 𝑎𝑛 is called a series.
To distinguish these two series, the latter is called a finite series. Series means an
infinite series. If the sequence of partial sums 𝑠𝑛 converges to s, we say that the series
∞
𝑎𝑛 converges or is convergent to the sum 𝑠 and write 𝑛=1 𝑎𝑛 = 𝑠 𝑜𝑟 𝑎𝑛 = 𝑠. If the
sequence 𝑠𝑛 is not convergent, we say that the series diverges or divergent. In
particular, if 𝑠𝑛 is divergent to +∞(−∞) or oscillating, we say that the series is
properly divergent to +∞(−∞) or oscillating, respectively.

CONVERGENT INTEGRAL: A improper integral whose limit exists and is finite.

CONVERGENT ITERATION: An iterative process for which the sequence of each state (in
order) is a convergent sequence

CONVERGENT NET: A net 𝑥𝑖 ; 𝑖 ∇ 𝐼 is said to be convergent to 𝑥 ∇ 𝑋 if for each

neighborhood 𝑈 of 𝑥 there is an index 𝑖 ∇ 𝐼 such that if 𝑗 ≥ 𝑖 then 𝑥𝑗 belongs to 𝑈. The
point 𝑥 is called a limit of the net 𝑥𝑖 ; 𝑖 ∇ 𝐼 and we often write 𝑥𝑖 → 𝑥. Convergence of
nets with index set 𝐼 = 𝑁 with the usual order is exactly convergence of sequences.

CONVERGENT PRODUCT: A continued product whose sequence of partial

products (analogues to partial sums) converges.

CONVERGENT SEQUENCE: It is a sequence in which the difference between two

consecutive terms tends to zero when n is large.
CONVERGENT SERIES: A series in which the sum is finite.
CONVERGES UNIFORMLY: Let 𝑋 be a set, (𝑌, 𝑑) a metric space and {𝑓𝑛 } a sequence of
functions from 𝑋 to 𝑌 , and 𝑓 ∶ 𝑋 → 𝑌 another function. If for any 𝜖 > 0 there exists an
integer 𝑁 such that
𝑑(𝑓𝑛 (𝑥), 𝑓(𝑥)) < 𝜖
for all 𝑛 > 𝑁 we say that fn converges uniformly to 𝑓.

CONVERSE: The converse of an implication 𝑝 ⇒ 𝑞 is the implication 𝑞 ⇒ 𝑝. If an

implication is true, then its converse may or may not be true.
CONVERSE OF THE PYTHAGOREAN THEOREM: If the square of the hypotenuse is equal
to the sum of the squares of the other two legs of a triangle, then the triangle is a right
triangle.
CONVERSION TABLE: A table showing the equivalent numerical values in two or more
desired units. When the value in one unit is given, the corresponding value in another
unit can be directly read from the table. One example is to convert from imperial to
metric measures and vice versa.
CONVEX (MANIFOLD): A subset K of a Riemannian manifold M is called convex if for any
two points in K there is a shortest path connecting them which lies entirely in K.

CONVEX COMBINATION: A convex combination of a finite number of points

𝑥1 , 𝑥2 , 𝑥3 , … , 𝑥𝑛 is defined as a point . 𝑥 = 𝜆1 𝑥1 + 𝜆2 𝑥2 = ⋯ 𝜆𝑛 𝑥𝑛 where 𝜆𝑖 is real and
𝑛
≥ 0, ∀ 𝑖 and 𝑖=1 𝜆𝑖 = 1.

CONVEX CONES: Let 𝑚 be a positive integer. A subset 𝐶 of 𝑅 𝑚 is said to be a convex cone

in 𝑅 𝑚 if 𝜆𝑣 + µ𝑤 ∇ 𝐶 for all 𝑣, 𝑤 ∇ 𝐶 and for all real numbers 𝜆 and µ satisfying
𝜆 ≥ 0 and µ ≥ 0. Every convex cone in 𝑅 𝑚 is a convex subset of 𝑅 𝑚 .

CONVEX FUNCTION: A function 𝑓 𝑥 is said to be strictly convex as x if for any two

other distinct points 𝑥1 and 𝑥2 , 𝑓 𝜆𝑥1 + 1 − 𝜆 𝑥2 } < 𝜆𝑓 𝑥1 + 1 − 𝜆 𝑓(𝑥2 ) , where
0 < 𝜆 < 1. On the other hand, a function 𝑓 𝑥 is strictly concave if -𝑓 𝑥 is strictly
convex.

CONVEX FUNCTION (DIFFERENTIAL GEOMETRY): A function 𝑓 on a Riemannian

manifold is a convex if for any geodesic 𝛾, the function 𝑓 ∘ 𝛾 is convex. A function 𝑓 is
called 𝜆-convex if for any geodesic 𝛾 with natural parameter 𝑡, the function 𝑓 ∘ 𝛾 𝑡 −
𝜆𝑡 2 is convex.

CONVEX HULL: The convex hull C(x) of any given set of points X is the set of all convex
combination of sets of points from X. For examples, if x is just the eight vertices of a
cube, then the convex hull C(x) is the whole cube.

CONVEX POLYHEDRON: The set of all convex combinations of finite number of points is
called the convex polyhedron generated by these points.

For example: The set of the area of the triangle is a convex polyhedron its vertices.

CONVEX POLYTOPE: A subset 𝑋 is said to be a convex polytope if there exist linear

functionals 𝜂1 , 𝜂2 , . . . , 𝜂𝑚 on 𝑅 𝑛 and real numbers 𝑠1 , 𝑠2 , . . . , 𝑠𝑚 such that 𝑋 = {𝑥 ∇ 𝑅 𝑛 ∶
𝜂𝑖 (𝑥) ≥ 𝑠𝑖 for 𝑖 = 1, 2, . . . , 𝑚}. Let (𝜂𝑖 ∶ 𝑖 ∇ 𝐼) be a finite collection of linear
functionals on 𝑅 𝑛 indexed by a finite set 𝐼, let 𝑠𝑖 be a real number for all 𝑖 ∇ 𝐼, and let
𝑋 = ⋂𝑖∇𝐼 {𝑥 ∇ 𝑅 ∶ 𝜂𝑖 (𝑥) ≥ 𝑠𝑖 }. Then 𝑋 is a convex polytope in 𝑅 𝑛 . A point 𝑥 of 𝑅 𝑛
belongs to the convex polytope 𝑋 if and only if 𝜂𝑖 (𝑥) ≥ 𝑠𝑖 for all 𝑖 ∇ 𝐼.
CONVEX SET: A plane or solid figure, such as a polygon or polyhedron, is said to be
convex if the line segment joining any two points inside it lies wholly inside it.

CONVEX FUZZY SET: A fuzzy set 𝑨 in 𝑹𝒏 is convex if and only if

𝜇𝐴 (𝜆𝑥1 + (1 − 𝜆𝑥2 ) ≥ min⁡
(𝜇𝐴 𝑥1 , 𝜇𝐴 𝑥2 )
for all 𝑥1 , 𝑥2 ∇ 𝑹𝒏 and all 𝜆 ∇ [0, 11.
CONVOLUTION THEOREM FOR FOURIER TRANSFORMS: Fourier transform of
convolution of 𝑓 𝑡 𝑎𝑛𝑑 𝑔 𝑡 is the product of their Fourier transforms i.e.

𝐹 𝐹∗𝑔 = 𝐹 𝑓 𝑡 𝐹 𝑔 𝑡

COORDINATE: One within a set of such numbers, called coordinates, which specifies the
position of a point e.g. abscissa and ordinate.

COORDINATE GEOMETRY: A study of geometric objects through their representation in

a coordinate system.

COORDINATE MAP OF A CHART: Let S be a manifold in 𝑅 𝑛 , and let 𝜍: 𝑈 → 𝑆 be a chart.

If 𝑝 ∇ 𝑆 we call (𝑥1 , . . . , 𝑥𝑚 ) ∇ 𝑈 the coordinates of 𝑝 with respect to 𝜍 when 𝑝 =
𝜍(𝑥). The map 𝜍 −1 : 𝜍(𝑈) → 𝑈 is called the coordinate map of 𝜍.
COORDINATE PLANE: A plane with a point selected as an origin, some length selected as
a unit of distance, and two perpendicular lines that intersect at the origin, with positive
and negative direction selected on each line. Traditionally, the lines are called 𝑥 (drawn
from left to right, with positive direction to the right of the origin) and 𝑦 (drawn from
bottom to top, with positive direction upward of the origin). Coordinates of a point are
determined by the distance of this point from the lines, and the signs of the coordinates
are determined by whether the point is in the positive or in the negative direction from
the origin.
COPLANAR SET OF POINTS: A set of points is coplanar if they all lie in the same plane.
Any three points are always coplanar. The vertices of a triangle are coplanar, but not the
vertices of a pyramid. Two lines are coplanar if they lie in the same plane, that is, if they
either intersect or are parallel.
COPRIME INTEGERS: Two integers n and m with no common factors are said to be
mutually prime or coprime. By definition, 𝒈𝒄𝒅(𝑛, 𝑚) = 1.
CORE: The core of a game with characteristic function v is the set, 𝐶(𝑣), of all
imputations that are not dominated for any coalition. The idea is that only such
imputations can persist in pre-game negotiations.
COROLLARY: Corollary is a result that follows from a theorem almost immediately, often
without further proof.
CORRELATION: Between two random variables, the correlation is a measure of the
extent to which a change in one tends to correspond to a change in the other. The
correlation is high or low depending on whether the relationship between the two is
close or not. If the change in one corresponds to a change in the other in the same
direction, there is positive correlation, and there is a negative correlation if the changes
are in opposite directions. Independent random variables have zero correlation. One
measure of correlation between the random variables 𝑋 and 𝑌 is the correlation
coefficient 𝜌 defined by
𝒄𝒐𝒗 (𝑥, 𝑦)
𝜌=
𝒗𝒂𝒓 𝑥 (𝒗𝒂𝒓 𝑦)
This satisfies — 1 ≤ 𝜌 ≤ 1. If 𝑋 and 𝑌 are linearly related, then 𝜌 = — 1 𝑜𝑟 + 1.
cos: It is the abbreviation for cosine function.
cosec: It is the abbreviation for cosecant function. (Also written as csc).
cosech: It is the abbreviation for hyperbolic cosecant function. (Also written as csch).
COSET: If 𝐻 is a subgroup of a group 𝐺, then for any element, 𝑥 of 𝐺 there is a left coset
𝑥𝐻 consisting of all the elements 𝑥𝑕, where 𝑕 is an element of 𝐻. Similarly, there is a
right coset, 𝐻𝑥, with elements 𝑕𝑥. If 𝑥𝐻 = 𝐻𝑥 ∀ 𝑥 ∇ 𝐺, then 𝐻 becomes a normal
subgroup of 𝐺.
COSET (CODING THEORY): Let 𝐶 be a linear code of length 𝑛 over 𝑓𝑄 and let ∇ 𝑓𝑄𝑁 .
Then the coset of 𝐶 determined by 𝑢 is defined to be the set 𝐶 + 𝑢 = {𝑐 + 𝑢 | 𝑐 ∇ 𝐶}.
For example, let 𝐶 = {000, 101, 010, 111} be a binary linear code. Then, 𝐶 + 000 =
𝐶, 𝐶 + 010 = {010, 111, 000, 101} = 𝐶, and 𝐶 + 001 = {001, 100, 011, 110}.
COSETS OF IDEAL IN A RING: Let 𝑅 be a ring and let 𝐼 be an ideal of 𝑅. The cosets of 𝐼 in
𝑅 are the subsets of 𝑅 that are of the form 𝐼 + 𝑥 for some 𝑥 ∇ 𝑅, where 𝐼 + 𝑥 = {𝑎 +
𝑥 ∶ 𝑎 ∇ 𝐼}.
We denote by 𝑅/𝐼 the set of cosets of 𝐼 in 𝑅. Let 𝑥 and 𝑥 ′ be elements of 𝑅. Then
𝐼 + 𝑥 = 𝐼 + 𝑥 ′ if and only if 𝑥 − 𝑥 ′ ∇ 𝐼.
Indeed if 𝐼 + 𝑥 = 𝐼 + 𝑥 ′ , then 𝑥 = 𝑐 + 𝑥 ′ for some 𝑐 ∇ 𝐼. But then 𝑥 − 𝑥 ′ = 𝑐, and
thus 𝑥 − 𝑥 ′ ∇ 𝐼.
Conversely if 𝑥 − 𝑥 ′ ∇ 𝐼 then 𝑥 − 𝑥 ′ = 𝑐 for some 𝑐 ∇ 𝐼. But then 𝐼 + 𝑥 = {𝑎 + 𝑥 ∶
𝑎 ∇ 𝐼} = {𝑎 + 𝑐 + 𝑥 ′ ∶ 𝑎 ∇ 𝐼} = {𝑏 + 𝑥 ′ ∶ 𝑏 ∇ 𝐼} = 𝐼 + 𝑥 ′ .
If 𝑥, 𝑥 ′ , 𝑦 and 𝑦 ′ are elements of 𝑅 satisfying 𝐼 + 𝑥 = 𝐼 + 𝑥 ′ and 𝐼 + 𝑦 = 𝐼 + 𝑦 ′
then
(𝑥 + 𝑦) − (𝑥 ′ + 𝑦 ′ ) = (𝑥 − 𝑥 ′ ) + (𝑦 − 𝑦 ′ ), 𝑥𝑦 − 𝑥 ′ 𝑦 ′ = 𝑥𝑦 − 𝑥𝑦 ′ + 𝑥𝑦 ′ −
𝑥 ′ 𝑦 ′ = 𝑥(𝑦 − 𝑦 ′ ) + (𝑥 − 𝑥 ′ )𝑦 ′ . But 𝑥 − 𝑥 ′ ∇ 𝐼 and 𝑦 − 𝑦 ′ ∇ 𝐼, and therefore
𝑥(𝑦 − 𝑦 ′ ) ∇ 𝐼 and (𝑥 − 𝑥 ′ )𝑦 ′ ∇ 𝐼, because 𝐼 is an ideal.
It follows that (𝑥 + 𝑦) − (𝑥 ′ + 𝑦 ′ ) ∇ 𝐼 and 𝑥𝑦 − 𝑥 ′ 𝑦 ′ ∇ 𝐼, and therefore 𝐼 + 𝑥 +
𝑦 = 𝐼 + 𝑥 ′ + 𝑦 ′ and 𝐼 + 𝑥𝑦 = 𝐼 + 𝑥 ′ 𝑦 ′ .
This shows that the quotient group 𝑅/𝐼 admits well-defined operations of addition and
multiplication, defined such that (𝐼 + 𝑥) + (𝐼 + 𝑦) = 𝐼 + 𝑥 + 𝑦 and (𝐼 + 𝑥)(𝐼 +
𝑦) = 𝐼 + 𝑥𝑦 for all 𝑥, 𝑦 ∇ 𝑅.
cosh: It is the abbreviation for hyperbolic cosine function.
COSINE RULE: In trigonometry, (also known as the cosine formula or the law
of cosines ) it is a statement about a general triangle that relates the lengths of its sides
to the cosine of one of its angles.

cot: It is the abbreviation for cotangent function.

COTANGENT SPACE: Let 𝑀 be an 𝑚-dimensional abstract manifold, and let 𝑝 ∇ 𝑀. The
dual space (𝑇𝑝 𝑀)∗ of 𝑇𝑝 𝑀 is denoted 𝑇𝑝∗ 𝑀 and called the cotangent space of 𝑀 at 𝑝. Its
elements are called cotangent vectors or just covectors.
COTERMINAL ANGLES: Angles in standard position with the same initial arm and
terminal arm.
coth: It is the abbreviation for hyperbolic cotangent function.
COUNTABLE ATLAS: An atlas of an abstract manifold M is said to be countable if the set
of charts in the atlas is countable. Let M be an abstract manifold. Then M has a countable
atlas if and only if there exists a countable base for the topology.
COUNTABLE SET: A set X is said to be a countable set if there is a one-to-one
correspondence between X and a subset of the set of natural numbers. Thus a countable
set is either finite or denumerable. Some authors use ‘countable’ to mean denumerable.
COUNTABLY ADDITIVE MEASURE: We say that µ is countably additive (or 𝜍 −additive)
if for any countable family (𝐴𝑛 ) of pairwise disjoint sets from 𝑅 we have µ(⋃𝑛 𝐴𝑛 ) =
𝑛 µ(𝐴𝑛 ). If the sum diverges, then as it will be the sum of positive numbers, we can,
without problem, define it to be +∞.

COUNTEREXAMPLE: An example that proves a statement false.

cov: It is the abbreviation for covariance of a pair of random variables.
COVARIANCE: A measure of the extent to which two variables vary together. When
divided by the product of the standard deviations of the two variables, this measure
becomes the Pearson (or product-moment) correlation coefficient. The covariance of
two random variables 𝑋 and 𝑌, denoted by 𝐶𝑜𝑣(𝑋, 𝑌), is equal to 𝐸((𝑋 — 𝜇𝑋 )(𝑌 — 𝜇𝑌 )),
where 𝜇𝑋 and 𝜇𝑌 are the population means of 𝑋 and 𝑌 respectively. If 𝑋 and 𝑌 are
independent random variables, then 𝐶𝑜𝑣(𝑋, 𝑌) = 0. For computational purposes, note
that 𝐸((𝑋 — 𝜇𝑋 )(𝑌 — 𝜇𝑌 )) = 𝐸(𝑋𝑌) — 𝜇𝑋 𝜇𝑌 . For a sample of n paired observations
(𝑥1 , 𝑦1 ), (𝑥2 , 𝑦2 ), … , (𝑥𝑛 , 𝑦𝑛 ), the sample covariance is equal to
𝑥𝑖 − 𝑥 𝑦𝑖 − 𝑦
𝑛
COVARIANT CONSTANTS: The covariant derivatives of the tensors 𝑔𝑖𝑗 , 𝑔𝑖𝑗 , 𝑔𝑗𝑖 vanish
identically. It means these tensors behave as constants in covariant differentiation.
Hence these tensors are defined a covariant constant.

COVARIANT DERIVATIVE OF A SCALAR: Covariant derivative of a scalar ∅ w.r.t. 𝑥 𝑖 is

𝜕∅
defined as its ordinary partial derivative w.r.t. 𝑥 𝑖 and is denoted by ∅, 𝑖. Thus ∅, 𝑖 = 𝜕𝑥 𝑖 .
𝜕∅
Evidently ∅, 𝑖 = 𝜕𝑥 𝑖 = ∆∅.

CRAMER, GABRIEL (1704–52): Gabriel Cramer was a Swiss mathematician whose

introduction to algebraic curves, published in 1750, contains the so-called Cramer’s rule
of solving a system of 𝑛 linear equations in 𝑛 variables. The rule was known earlier by
Maclaurin.
CRAMER’S RULE: Consider a set of 𝑛 linear equations in 𝑛 unknowns 𝑥1 , 𝑥2 , … , 𝑥𝑛 ,
written in matrix form as 𝑨𝒙 = 𝒃. When 𝑨 is an invertible matrix, the set of equations
has a unique solution 𝒙 = 𝑨—1 𝒃. Since 𝑨—1 = (1/𝑑𝑒𝑡 𝑨) 𝑎𝑑𝑗 𝑨, this gives the solution
𝑎𝑑𝑗𝐴 𝑏
𝑥=
det 𝐴
The equation ax + by + cz = 0 is a linear homogeneous equation in x, y and z. On the
other hand the equation ax + by + cz = 0 is a linear non-homogeneous equation in
x, y and z.

Let the 𝑛 linear simultaneous equations in 𝑛 unknowns 𝐴, x2 , x3 , … . , xn be

a11 x1 − a12 x2 + ⋯ + a1n xn = b1 ,

a21 x1 − a22 x2 + ⋯ + a2n xn = b2 ,

… … … … … … … …

… … … … … … … …

an1 x1 − an2 x2 + ⋯ + ann xn = bn .

a11 a12 ⋯ a1n

a21 a22 ⋯ a2n
Let ∆= ⋯ ⋯ ⋯ ⋯ ≠ 0.
⋯ ⋯ ⋯ ⋯
an1 an2 ⋯ ann

Suppose A11 , A12 , A13 ,………. etc. denote the cofactors if a11 , a12 , a13, ………. etc, in ∆. Then
multiplying the given equation respectively byA11 , A12 , A13 ,……… An1 and adding, we get

𝑥1 a11 A11 + a21 A21 + ⋯ + an1 An1 + 𝑥2 0 + 𝑥2 0 + 𝑥3 0 + … + 𝑥𝑛 0

= b1 A11 + b2 A21 + ⋯ + bn An1

Or 𝑥1 ∆= ∆1 ,

Where ∆1 is the determinant obtained by replacing the elements in the first column of ∆
by the elements b1 , b2 , … , bn .
Again multiplying the given equation respectively by A12 , A22 , … , An2 and adding, we get

𝑥2 ∆= ∆2 ,

Where ∆2 is the determinant obtained by replacing the elements in the second column
of ∆ by the elements b1 , b2 , … , bn .

Similarly, we get

𝑥3 ∆= ∆3

…………

𝑥𝑛 ∆= ∆𝑛 .

This method of solving 𝑛 simultaneous equations in 𝑛 unknowns is known as Cramer’s

Rule.

Thus by Cramer’s rule, if ∆≠ 0, we have

∆i
xi = , i = 1,2, … , n,
∆

Where ∆𝑖 is the determinant obtained by replacing the ith column of ∆ by the elements
b1 , b2 , … , bn .

CRITICAL POINT: A critical point for a function is a point where the first derivative(s) is
(are) zero.
CRITICAL POINTS (COMPLEX ANALYSIS): Consider the bilinear transformation

𝑎𝑧 + 𝑏
𝑤=𝑇 𝑧 =
𝑐𝑧 + 𝑑

Solving this for 𝑧, we get the inverse map as

𝑏 − 𝑤𝑑
𝑧 = 𝑇 −1 𝑤 =
𝑤𝑐 − 𝑎

The transformation 𝑇 associates a unique point of the 𝑤 − 𝑝𝑙𝑎𝑛𝑒 to any point of

𝑑
𝑧 − 𝑝𝑙𝑎𝑛𝑒 except the point 𝑧 = − 𝑐 when 𝑐 ≠ 0. The transformation 𝑇 −1 associates a
𝑎
unique point of 𝑧 − 𝑝𝑙𝑎𝑛𝑒 to any point 𝑤 plane except the point 𝑤 = , when 𝑐 ≠ 0.
𝑐
 𝑑 𝑎
These exceptional points 𝑧 = − 𝑐 and 𝑤 − 𝑐 are mapped into the points 𝑤 = ∞ 𝑎𝑛𝑑 𝑧 =

∞ respectively are obvious from (1) and (2)

𝑑𝑤 𝑎𝑑 −𝑏𝑐
From (1) = (𝑐𝑧 +𝑑)²
𝑑𝑧

𝑑
𝑑𝑤 ∞ 𝑖𝑓 𝑧 = − 𝑐
This ⟹ =
𝑑𝑧
0 𝑖𝑓 𝑧 = ∞

𝑑
The points 𝑧 = − 𝑐 , 𝑧 = ∞ are called Critical points where the conformal property does

not hold good.

CRITICAL REGION: If the calculated value of a test statistic falls within the critical region,
then the null hypothesis is rejected.
CROSS RATIO: For a set of coplanar points 𝐴, 𝐵, 𝐶, 𝐷, it is,
𝐴𝐶 × 𝐵𝐷
𝐴𝐷 × 𝐵𝐶
and for points 𝑧1 , 𝑧2 , 𝑧3 , 𝑧4 in the complex plane it is
𝑧1 − 𝑧3 𝑧2 − 𝑧4
𝑧1 − 𝑧4 𝑧2 − 𝑧3
The definition in the complex plane can be extended to the Riemann sphere by
continuity.
CRYPTOGRAPHY: Cryptography is the area of mathematics concerning the secure
coding and decoding of information, often relying on mathematics such as prime
factorizations of very large numbers.
csc: It is the abbreviation for cosecant function. (Also written as cosec).
csch: It is the abbreviation for hyperbolic cosecant function. (Also written as cosech).
CUBE ROOT OF UNITY: A complex number z such that 𝑧 3 = 1. The three cube roots of
−1+𝑖 3 −1−𝑖 3
unity are 1, 𝜔 and 𝜔2 , where 𝜔 = and 𝜔2 = .
2 2

CUBIC EQUATION: A polynomial having a degree of 3 (i.e. the highest power is 3), of the
form 𝑎𝑥 3 + 𝑏𝑥 2 + 𝑐𝑥 + 𝑑 = 0, which can be solved by factorization or formula to
find its three roots.
CUMULATIVE DISTRIBUTION FUNCTION: For a random variable 𝑋, the cumulative
distribution function (c.d.f.) is the function 𝐹 defined by 𝐹(𝑥) = 𝑃(𝑋 ≤ 𝑥). Thus, for a
discrete random variable
 𝐹 𝑥 − 𝑝(𝑥𝑖 )
𝑥 𝑖 ≤𝑋

where p is the probability mass function, and, for a continuous random variable,
𝑥

𝐹 𝑥 = 𝑓 𝑡 𝑑𝑡
−∞

where f is the probability density function.

CUMULATIVE FREQUENCY: The sum of the frequencies of all the values up to a given
value. If the values 𝑥1 , 𝑥2 , … , 𝑥𝑛 , in ascending order, occur with frequencies 𝑓1 , 𝑓2 , … , 𝑓𝑛 ,
respectively, then the cumulative frequency at 𝑥𝑖 is equal to 𝑓1 + 𝑓2 , … + 𝑓𝑖 . Cumulative
frequencies may be similarly obtained for grouped data.
CURL OF A VECTOR: For a vector function of position 𝑽 (𝒓) = 𝑉𝑥 𝒊 + 𝑉𝑦 𝒋 + 𝑉𝑧 𝒌, the curl
𝜕 𝜕 𝜕
of 𝑽 is the vector product of the operator del ∆= 𝒊 𝜕𝑥 + 𝒋 𝜕𝑦 + 𝒌 𝜕𝑧 , with 𝑽 giving

𝜕𝑽 𝜕𝑽 𝜕𝑽
𝑐𝑢𝑟𝑙 𝑽 = ∆ × 𝑽 = 𝒊 × +𝒋× +𝒌×
𝜕𝑥 𝜕𝑦 𝜕𝑧
which can be written in determinant form as
𝑖 𝑗 𝑘
𝜕 𝜕 𝜕
𝜕𝑥 𝜕𝑦 𝜕𝑧
𝑉𝑥 𝑉𝑦 𝑉𝑧
CURVATURE: The rate of change of direction of a curve at a point 𝑃 on the curve. The
Greek letter 𝜅 is used to denote curvature and
𝑦 ′′
𝜅= .
1 + 𝑦′ 2 3/2

1
𝜌 = 𝜅 is the radius of curvature which is the radius of the circle which best fits the curve

at that point, matching the position, the gradient and the second differential of the curve
at that point. The centre of curvature is the centre of that best-fitting circle, known as
the circle of curvature.
CURVE: A curve is a continuous mapping of the segment [0, 1] into another space -
container of the curve. A curve may not look as a line. For example, there are space
filling curves.
CUSP: A point at which two or more branches of a curve meet, and at which the limits of
the tangents approaching that point along each branch coincide. There are two main
characteristics used to describe cusps. In a single or simple cusp there are only two
branches, and the limits of the second differentials approaching that point are different.
If the branches are on opposite sides of the common tangent, it is said to be a cusp of the
first kind, and if the branches are on the same side of the common tangent it is a cusp of
the second kind.

A double cusp or point of osculation has four branches, comprising two continuously
differentiable curves meeting at a point with a common tangent. Double cusps can also
be of the first or second kind, or one or both curves can have a point of inflexion at the
cusp, so the tangent intersects the curve, in which case it is a point of osculinflection.

C(X) SPACES: All our topological spaces are assumed Hausdorff. Let X be a compact
space, and let 𝐶𝐾 (𝑋) be the space of continuous functions from 𝑋 to 𝐾, with pointwise
operations, so that 𝐶𝐾 (𝑋) is a vector space. We norm 𝐶𝐾 (𝑋)by setting 𝑓 ∞ =
𝑠𝑢𝑝
𝑓(𝑥) ; 𝑓 ∇ 𝐶𝐾 (𝑋).
𝑥∇𝑋

Note that
  Let 𝑋 be a compact space. Then 𝐶𝐾 (𝑋) is a Banach space.
 Let 𝐸 be a vector space, and let || · || (1) and || · || (2) be norms on 𝐸. These norms
are equivalent if there exists 𝑚 > 0 with 𝑚−1 𝑥 (2) ≤ 𝑥 (1) ≤𝑚 𝑥 (2) . (𝑥 ∇
𝐸)
 Let E be a finite-dimensional vector space with basis {𝑒1 , … , 𝑒𝑛 }, so we can
identify 𝐸 with 𝐾 𝑛 as vector spaces, and hence talk about the norm || · || (2) on E.
If || · || is any norm on 𝐸, then || · || and || · || (2) are equivalent.
 Let E be a finite-dimensional normed space. Then a subset 𝑋 ⊆ 𝐸 is compact if
and only if it is closed and bounded.
 Let 𝐸 be a normed vector space, and let 𝐹 be a closed subspace of 𝐸 with 𝐸 ≠ 𝐹.
For 0 < 𝜃 < 1, we can find 𝑥0 ∇ 𝐸 with ||𝑥0 || ≤ 1 and ||𝑥0 − 𝑦|| > 𝜃 for 𝑦 ∇ 𝐹.
 Let 𝐸 be an infinite-dimensional normed vector space. Then the closed unit ball
of 𝐸, the set {𝑥 ∇ 𝐸 ∶ ||𝑥|| ≤ 1}, is not compact.

CYCLE (GRAPH THEORY): A closed path with at least one edge. In a graph, a cycle is a
sequence 𝑣0 , 𝑒1 , 𝑣1 , … , 𝑒𝑘 , 𝑣𝑘 (𝑘 ≥ 1) of alternately vertices and edges (where 𝑒𝑖 is an
edge joining 𝑣𝑖—1 and 𝑣𝑖 ), with all the edges different and all the vertices different,

except that 𝑣0 = 𝑣𝑘 .
CYCLIC GROUP: Let 𝑎 be an element of a multiplicative group 𝐺. The elements 𝑎𝑟 , where
𝑟 is an integer (positive, zero or negative), form a subgroup of 𝐺, called the subgroup
generated by 𝑎. A group 𝐺 is cyclic if there is an element a in 𝐺 such that the subgroup
generated by 𝑎 is the whole of 𝐺. If 𝐺 is a finite cyclic group with identity element 𝑒, the
set of elements of 𝐺 may be written as {𝑒, 𝑎, 𝑎2 , … , 𝑎𝑛−1 }, where 𝑎𝑛 = 𝑒 and 𝑛 is the
smallest such positive integer. If 𝐺 is an infinite cyclic group, the set of elements may be
written { … , 𝑎−2 , 𝑎−1 , 𝑒, 𝑎, 𝑎2 , … }.
CYCLIC POLYGON: A polygon whose vertices all lie on the same circle. All triangles, all
rectangles, and all regular polygons are cyclic. Convex quadrilaterals, whose opposite
angles are supplementary, are also cyclic.
CYCLIC QUADRILATERAL: A four-sided figure whose four vertices all lie on a
circumscribed circle. If the opposite angles of a quadrilateral add to 180°, it is a cyclic
quadrilateral.
CYCLOID: The curve traced out by a point on the circumference of a circle that rolls
without slipping along a straight line. With suitable axes, the cycloid has parametric
equations 𝑥 = 𝑎(𝑡 — 𝑠𝑖𝑛 𝑡), 𝑦 = 𝑎(1 — 𝑐𝑜𝑠 𝑡)(𝑡 ∇ 𝑹), where 𝑎 is a constant (equal
to the radius of the rolling circle). In the figure, 𝑂𝐴 = 2𝜋𝑎.

CYCLOTOMIC NUMBER FIELDS: The starting point for cyclotomic number fields is the
irreducibility of the cyclotomic polynomial 𝑓(𝑡) = 𝑡 𝑛 −1 + 𝑡 𝑛 −2 + … + 𝑡 + 1 in 𝑄[𝑡]
for any positive rational prime 𝑝. To see this, note that 𝑓(𝑡 + 1) = 𝑡 𝑝 −1 +
𝑝(𝑡 𝑝−2 + . . . ) + 𝑝, and the irreducibility now follows immediately from Eisenstein’s
criterion. For any positive rational prime 𝑝, the 𝑝-th cyclotomic number field is the
algebraic number field 𝑄(𝜁), where 𝜁 = 𝑒 2𝜋𝑖 /𝑝 is a primitive 𝑝-th root of unity.
CYCLOTOMIC POLYNOMIAL Φn(z): The polynomial whose roots are all the primitive n-
th roots of unity. We know
𝑧 𝑛 — 1 ≡ (𝑧 — 1)(𝑧 𝑛−1 + 𝑧 𝑛−2 + … + 𝑧 + 1),
so when n is prime 𝛷𝑛 (𝑧) = 𝑧 𝑛 −1 + 𝑧 𝑛 −2 + … + 𝑧 + 1 is the cyclotomic polynomial,
but when 𝑛 = 4, 𝑧 4 — 1 = (𝑧 — 1)(𝑧 + 1)(𝑧 2 + 1), and so 𝛷4 (𝑧) = (𝑧 2 + 1).
CYLINDRICAL HELICES: A cylindrical helix is a curve traced out on the surface of the
cylinder and cuts the generators at a constant angle.
Thus the tangent at any point on the helix makes a constant angle say 𝛼, with a fixed
line, this fixed line is known as axis (generator) of the cylinder, the axis of the cylinder is
also called the axis of the helix.
CYLINDRICAL POLAR COORDINATES: Suppose that three mutually perpendicular
directed lines 𝑂𝑥, 𝑂𝑦 and 𝑂𝑧, intersecting at the point 𝑂 and forming a right-handed
system, are taken as coordinate axes. For any point 𝑃, let 𝑀 and 𝑁 be the projections of
𝑃 onto the 𝑥𝑦 −plane and the 𝑧 −axis respectively. Then 𝑂𝑁 = 𝑃𝑀 = 𝑧, the
𝑧 −coordinate of 𝑃. Let 𝜌 = |𝑃𝑁|, the distance of 𝑃 from the 𝑧 −axis, and let ø be the
angle ∠𝑥𝑂𝑀 in radians (0 ≤ ø < 2 𝜋). Then
(𝜌, ø, 𝑧) are the cylindrical polar coordinates of 𝑃.
(It should be noted that the points of the 𝑧 −axis
give no value for ø.) The two coordinates (𝜌, ø) can be seen as polar coordinates of the
point 𝑀 and, as with polar coordinates, ø + 2 𝑘𝜋, where 𝑘 is an integer, may be allowed
in place of ø.

D
D’ ALEMBERT’S PARADOX: If we super- impose on the system a uniform velocity 𝑈 in
the direction opposite to that of the current the inviscid fluid at a great distance remain
undisturbed and the body moves with uniform velocity 𝑈. The dynamical conditions
remain unaltered by superposing a uniform velocity therefore, the resistance to a body
moving with uniform velocity through an unbounded inviscid fluid, otherwise at rest,
vanishes. This is in contradiction to the common experience and is called 𝐷’ Alembert’s
Paradox.

D’ ALEMBERT’S ROOT TEST: Let 𝑢𝑛 be a series of positive terms such that

𝑢 𝑛 +1
lim𝑛→∞ = 𝑙. Then
𝑢𝑛

1. 𝑢𝑛 converges if 𝑙 < 1.
2. 𝑢𝑛 diverges if 𝑙 > 1.
3. The test fails if 𝑙 = 1.
DARBOUX FUNCTION: For topological spaces 𝑋 and 𝑌, a function 𝑓 ∶ 𝑋 → 𝑌 is a
Darboux function (or has the Darboux property), 𝑓 ∇ 𝐷(𝑋, 𝑌 ), provided the image
𝑓 [𝐶] of 𝐶 under 𝑓 is a connected subset of 𝑌 for every connected subset 𝐶 of 𝑋. In
particular, 𝑓 ∶ 𝑅 → 𝑅 is Darboux provided 𝑓 maps intervals onto intervals, that is,
when it has the intermediate value property.
DARBOUX THEOREM: Let 𝑓 be a bounded function defined on [𝑎, 𝑏]. Then to every ∇> 0,
there corresponds 𝛿 > 0 such that
𝑏

𝑈 𝑃, 𝑓 < 𝑓𝑑𝑥+∇
𝑎

and
𝑏

𝐿 𝑃, 𝑓 > 𝑓𝑑𝑥−∇
𝑎

DATA: The observations gathered from an experiment, survey or observational study.

Often the data are a randomly selected sample from an underlying population.
Numerical data are discrete if the underlying population is finite or countably infinite
and are continuous if the underlying population forms an interval, finite or infinite. Data
are nominal if the observations are not numerical or quantitative, but are descriptive
and have no natural order. Data specifying country of origin, type of vehicle or subject
studied, for example, are nominal. Note that the word ‘data’ is plural. The singular
‘datum’ may be used for a single observation.

DATA ANALYTICS: Data analytics (DA) is the science of examining raw data with the
purpose of drawing conclusions about that information. Data analytics is used in many
industries to allow companies and organization to make better business decisions and
in the sciences to verify or disprove existing models or theories. Data analytics is
distinguished from data mining by the scope, purpose and focus of the analysis. Data
miners sort through huge data sets using sophisticated software to identify
undiscovered patterns and establish hidden relationships. Data analytics focuses on
inference, the process of deriving a conclusion based solely on what is already known by
the researcher.

d-DIMENSIONAL WEIERSTRASS THEOREM: The restrictions of the polynomials in 𝑑

arguments to any compact subset 𝐾 of 𝑅 𝑑 is dense in 𝐶(𝐾).

DE BRANGES'S THEOREM: Also known as Bieberbach conjecture, is a theorem that gives

a necessary condition on a holomorphic function in order for it to map the open unit
disk of the complex plane injectively to the complex plane. It was posed by Ludwig
Bieberbach (1916) and finally proven by Louis de Branges (1985).
The statement concerns the Taylor coefficients 𝑎𝑛 of such a function, normalized as is
always possible so that 𝑎0 = 0 and 𝑎1 = 1. That is, we consider a function defined on
the open unit disk which is holomorphic and injective (univalent) with Taylor series of
the form

The theorem states that

DECISION ANALYSIS: The branch of mathematics considering strategies to be used

when decisions have to be made at stages in a process, but the outcomes resulting from
the implementation of those decisions are dependent on chance.
DECISION THEORY: The area of statistics and game theory concerned with decision
making under uncertainty to maximize expected utility.
DECISION TREE: The diagram used to represent the process in a decision analysis
problem. Different symbols are used to denote the different types of node or vertex. For
example, decisions may be shown as rectangles, chance events as circles, and payoffs as
triangles.
DECISION PROBLEMS: A decision problem (or recognition problem) is one that takes
the form of a question with a yes/no answer.
DECISION VARIABLES: The quantities to be found in linear programming or other
constrained optimization problems.
DECOMPOSITION THEOREM: If 𝐺 is a subspace (i.e. a closed linear manifold) of a Hilbert
space 𝐻, then every vector 𝑕 of H can be uniquely decomposed into the sum of a vector
𝑔 belonging to 𝐺 and a vector 𝑓 which is orthogonal to : i.e.

𝐻 = 𝐺 ⊕ 𝐺⊥

Morever, if 𝑑 = 𝑑𝑖𝑠𝑡𝑎𝑛𝑐𝑒 𝑜𝑓 𝑕 𝑓𝑟𝑜𝑚 𝐺

= 𝑖𝑛𝑓𝑖𝑚𝑢𝑚 𝑜𝑓 𝑕 − 𝑔′ ,

Taken over all the vectors 𝑔′ 𝑜𝑓 𝐺, 𝑡𝑕𝑒𝑛 𝑓 = 𝑕 − 𝑔 − 𝑑, 𝑖. 𝑒. this infimum is attained

at the vector 𝑔.
DECREASING FUNCTION: A function f (x) is a decreasing function if f (a) _ f (b) when a _
b.
DECREASING SEQUENCE: A real sequence 𝑎1 , 𝑎2 , 𝑎3 , … is said to be a decreasing
sequence if 𝑎𝑛 ≥ 𝑎𝑛+1 for all 𝑛 and strictly decreasing if 𝑎𝑛 > 𝑎𝑛+1 for all 𝑛.
DEDEKIND- CANTOR AXIOM: To every real number there corresponds a unique point on
a directed line and conversely, to every point on a directed line there corresponds a
unique real number.
DEDEKIND DOMAINS: A Dedekind domain is an integrally-closed Noetherian domain in
which every non-zero prime ideal is maximal. It follows from this definition that a unital
commutative ring R is a Dedekind domain if and only if it possesses all four of the
following properties:
(i) R is an integral domain;
(ii) every ideal of R is finitely generated,
(iii) R is integrally closed in its field of fractions;
(iv) every non-zero prime ideal of R is maximal.
Properties (i) and (ii) characterize Noetherian domains, and property (iii) characterizes
integrally-closed domains. Every principal ideal domain is a Dedekind domain.
DEDEKIND, JULIUS WILHELM RICHARD (1831–1916): Richard Dedekind was a German
mathematician who developed a formal construction of the real numbers from the
rational numbers by means of the so-called Dedekind cut. This new approach to
irrational numbers, contained in the very readable paper Continuity and Irrational
Numbers, was an important step towards the formalization of mathematics. He also
proposed a definition of infinite sets that was taken up by Cantor, with whom he
developed a lasting friendship.
DEDEKIND’S AXIOM: Let 𝐿 and 𝑈 be two subsets of 𝑅 such that
1. 𝐿 ≠ 𝜑, 𝑈 ≠ 𝜑
2. 𝐿 ∪ 𝑈 = 𝑅
3. 𝑥 ∇ 𝐿, 𝑦 ∇ 𝑈 ⇒ 𝑥 < 𝑦
Then the subset L has the greatest member or the subset U has the smallest member i.e.
∃𝛼 ∇ 𝑅 𝑠. 𝑡. 𝑥 < 𝛼 ⇒ 𝑥 ∇ 𝐿, 𝑦 > 𝛼 ⇒ 𝑦 ∇ 𝑈.
DEDUCTION: A deduction is a conclusion arrived at by reasoning.
DEDUCTIVE REASONING: Using facts, definitions, accepted properties, and the laws of
logic to make a logical argument.
DEFICIENT NUMBER: A positive integer is said to be a deficient number if it is larger
than the sum of its positive divisors
DEFINITE INTEGRAL: If 𝑓(𝑥) represents a function of 𝑥 that is always nonnegative, then
the definite integral of 𝑓(𝑥) between 𝑎 and 𝑏 represents the area under the curve
𝑦 = 𝑓 (𝑥), above the x-axis, to the right of the line 𝑥 = 𝑎, and to the left of the line x = b.
The definite integral is represented by the expression
𝑏

𝑓 𝑥 𝑑𝑥
𝑎

and 𝑎 and 𝑏 are the limits of integration.

DEFINITION: An exact statement of the meaning, nature, and/or limits of a
mathematical object.
deg: It is the abbreviation for degree of a polynomial. (Also written as ∂ .
DEGENERACY IN TRANSPORTATION PROBLEMS: The solution procedure for non-
degenerate basic feasible solution with exactly m+n-1 strictly positive allocations in
independent positions has been discussed so far. However, sometimes it is not possible
to get such initial feasible solution to start with. Thus degeneracy occurs in the
transportation problem whenever a number of occupied cells is less than m+n-1.

Basic feasible solution to an m –origin and n- destination transportation problem can

have at most m+n-1 number of positive (non –zero) basic variables. If this number is
exactly m+n-1 , the BFS is said to be non-degenerate; and if less than m+n-1 the basic
solution degenerates. It follows that whenever the number of basic cells is less than
m+n-1 , the transportation problem is a degenerate one.

Degeneracy in transportation problems can occur in two ways:

(i) Basic feasible solutions may be degenerate from the initial stage onward.
(ii) They may become degenerate at any intermediate stage.

A special technique called as the Assignment technique is used to solve the special type
of problems called Assignment problem. This type of problems may be defined as the
traditional or classical problems where the main motive is to allot a number of origins
to the equal number of destinations at a least cost are called assignment problems.
The assignment problems can be stated in the form of 𝑛 × 𝑛 matrix 𝑐𝑖𝑗 of real
numbers called cost matrix or effectiveness matrix.

DEGENERACY PROBLEM (TIE FOR MINIMUM RATIO): At the stage of improving the
solution during simplex procedure, minimum ratio 𝑥8 /𝑥𝑘 (𝑥𝑘 > 0) is determined in the
last column of simplex table to find the key row (i.e. a row containing the key element).
But, sometimes this ratio may not be unique, i.e, the key element (hence the variable to
leave the basis) is not uniquely determined or at the very first iteration, the value of one
or more basic variables in the XB column become equal to zero, this causes the problem
of degeneracy.

However, if the minimum ratio is zero for two or more basic variables, degeneracy may
result the simplex routine to cycle indefinitely. That is, the solution which we have
obtained in one iteration may repeat again after few iterations and therefore no
optimum solution may be obtained under such circumstances.

DEGENERATE CONIC SECTION: A point, line, or pair of lines that arise as a limiting form
of a conic.
DEGREE OF AN ALGEBRAIC NUMBER: The degree of an algebraic number field 𝐾 is the
dimension [𝐾: 𝑄] of K considered as a vector space over the field 𝑄 of rational numbers.
DEGREE OF A VERTEX OF A GRAPH: The degree of a vertex 𝑉 of a graph is the number of
edges ending at 𝑉. If loops are allowed, each loop joining 𝑉 to itself contributes two to
the degree of 𝑉.
DEGREES OF FREEDOM (MECHANICS): The number of degrees of freedom of a body is
the minimum number of independent coordinates required to describe the position of
the body at any instant, relative to a frame of reference. A particle in straight-line
motion or circular motion has one degree of freedom. A rigid body rotating about a fixed
axis also has one degree of freedom. A particle moving in a plane, such as a projectile, or
a particle moving on a cylindrical or spherical surface has two degrees of freedom. A
rigid body in general motion has six degrees of freedom.
DEGREES OF FREEDOM (STATISTICS): A positive integer normally equal to the number
of independent observations in a sample minus the number of population parameters to
be estimated from the sample. When the chi-squared test is applied to a contingency
table with 𝑕 rows and 𝑘 columns, the number of degrees of freedom equals
(𝑕 – 1)(𝑘 – 1).
DEL: The del symbol ∆ is used to represent this vector of differential operators:
𝜕 𝜕 𝜕
∆= , ,
𝜕𝑥 𝜕𝑦 𝜕𝑧
DELETED NEIGHBOURHOOD OF A POINT: If from a nbd of a point 𝑝, the point 𝑝 itself is
deleted, then we get a deleted nbd of the point 𝑝.
DELTA FUNCTION 𝛅 𝐱 − 𝐚 : It is defined as

∞ if x = a
δ(x − a) =
0 if x ≠ a
∞
together with the condition ∫−∞ δ x − a dx = 1.

DE MOIVRE, ABRAHAM (1667–1754): De Moivre was a prolific mathematician, born in

France, who later settled in England. In De Moivre’s Theorem, he is remembered for his
use of complex numbers in trigonometry.
DE MOIVRE’S THEOREM: For all positive integers 𝑛,
(𝑐𝑜𝑠 𝜃 + 𝑖 𝑠𝑖𝑛 𝜃)𝑛 = 𝑐𝑜𝑠 𝑛𝜃 + 𝑖 𝑠𝑖𝑛 𝑛𝜃.
DE MORGAN, AUGUSTUS (1806–71): Augustus De Morgan was a British mathematician
and logician who was responsible for developing a more symbolic approach to algebra,
and who played a considerable role in the beginnings of symbolic logic. His name is
remembered in De Morgan’s laws, which he formulated.
DE MORGAN’S LAWS: For all sets A and B (subsets of a universal set),
(𝐴 ∪ 𝐵)′ = 𝐴′ ∩ 𝐵′ 𝑎𝑛𝑑 (𝐴 ∩ 𝐵)′ = 𝐴′ ∪ 𝐵′.
DE MORGAN AND BERTRAND’S TEST: The series 𝑢𝑛 of positive terms is convergent or
divergent according as
𝑢𝑛
lim 𝑛 − 1 − 1 log 𝑛 > 1 𝑜𝑟 < 1.
𝑛 →∞ 𝑢𝑛 +1
DENSE IN ITSELF: A subset 𝐴 of 𝑅 is said to be dense in itself if it possesses no isolated
point i.e. every point of 𝐴 is a limit point of 𝐴.
DENSENESS PROPERTY OF REAL NUMBERS: Between any two distinct real numbers
there always lies a rational number and therefore infinitely many rational numbers.
DENSE SUBSET: A subset 𝑅 of a normed linear space 𝑁 is said to be dense in 𝑁, if every
element of 𝑁 is the limit of some fundamental sequence of vectors in 𝑅. This is
equivalent to saying that given any vector 𝑓 of 𝑁 and given any 𝜀 > 0, a vector 𝑔∇ in 𝑅
can be found such that 𝑓 − 𝑓∇ < 𝜀.

DENSITY PROPERTY: The property that states that there always exists another rational
number between any two given rational numbers. This means that the set of rational
numbers is dense.
DENSITY PROPERTY FOR REAL NUMBERS: The property that states that there always
exists another real number between any two given real numbers. This means that the
set of real numbers is dense.
DENUMERABLE SET: A set 𝑋 is said to be a denumerable set if there is a one-to-one
correspondence between X and the set of natural numbers. It can be shown that the set
of rational numbers is denumerable but that the set of real numbers is not.
DEPENDENT EQUATIONS: A set of equations where at least one of the set may be
expressed as a linear combination of the others.
DEPENDENT EVENT: An event that is affected by the occurrence of other events
DEPENDENT VARIABLE: The dependent variable stands for the output number of a
function. In the equation 𝑦 = 𝑓 (𝑥), 𝑦 is the dependent variable and 𝑥 is the
independent variable. The value of 𝑦 depends on the value of 𝑥. In statistics, the variable
which is thought might be influenced by certain other explanatory variables. In
regression, a relationship is sought between the dependent variable and the
explanatory variables. The purpose is normally to enable the value of the dependent
DERIVATIVE: The derivative of a function is the rate of change of that function. On the
graph of the curve 𝑦 = 𝑓 (𝑥), the derivative at 𝑥 is equal to the slope of the tangent line
at the point (𝑥, 𝑓 (𝑥)). If the function represents the position of an object as a function of
time, then the derivative represents the velocity of the object. Derivatives can be
calculated from this expression:
Function Derivative
𝑦 = 𝑓(𝑥) 𝑓 𝑥 + 𝑕 − 𝑓(𝑥)
𝑦 ′ = 𝑓 ′ 𝑥 = lim
𝑕→0 𝑕
DERIVED SET: The set of all limit points of 𝐴 is called derived set of 𝐴.
DERIVED UNIT: It is a unit that is defined or obtained in terms of the fundamental units.
For example, a unit of force, one Newton (N), is defined as the amount of force that will
accelerate a mass of one kilogram at the rate of one meter per second per second.
Therefore, its formal definition is kg·m·s-2.
DEROGATORY AND NON-DEROGATORY MATRIX: An 𝑛-rowed matrix is said to be
derogatory or non-derogatory, according as the degree of its minimal equation is less
than or equal to 𝑛.

If the roots of the characteristic equation of a matrix are all distinct, then the matrix is
non-derogatory.

DESCARTES, RENE: Rene Descartes (1596 to 1650) was a French mathematician and
philosopher who is noted for the sentence “I think, therefore I am” and for developing
the concept now known as rectangular, or Cartesian coordinates. He was the
mathematician who in mathematics is known mainly for his methods of applying
algebra to geometry, from which analytic geometry developed. He expounded these in
La Géométrie, in which he was also concerned to use geometry to solve algebraic
problems.
DESCARTES’ RULE OF SIGNS: Descartes’ rule of signs states that the number of positive
roots of a polynomial equation will equal the number of sign changes among the
coefficients, or that number minus a multiple of 2. To count the sign changes, be sure the
polynomial terms are arranged in descending order by power of x, and ignore any zero
coefficients.
DESCENDING CHAIN CONDITION: A collection 𝑆 of subsets of a set 𝑋 satisfies the
ascending chain condition or DCC if there does not exist an infinite descending chain
𝑆1 ⊇ 𝑆2 ⊇ · · · of subsets from 𝑆.
DESCRIPTIVE GEOMETRY: A method of representing three-dimensional objects by
projections on the two-dimensional plane using a specific set of procedures.
DESCRIPTIVE STATISTICS: This is the part of the subject of statistics concerned with
describing the basic statistical features of a set of observations. Simple numerical
summaries, using notions such as mean, range and standard deviation, together with
appropriate diagrams such as histograms, are used to present an overall impression of
the data. Descriptive statistics is the study of ways to summarize data. For example, the
mean, median, and standard deviation are descriptive statistics that summarize some of
the properties of a list of numbers.
det: It is the abbreviation for determinant of a matrix or linear transformation.
DETERMINANT: The determinant of a matrix is a number that is useful in describing the
characteristics of the matrix. The determinant is symbolized by enclosing the matrix in
vertical lines. The determinant of a 2 × 2 matrix is:
𝑎 𝑏
= 𝑎𝑑 − 𝑏𝑐
𝑐 𝑑
The determinant of a 3 × 3 matrix can be found from:
𝑎 𝑏 𝑐
𝑓 𝑑 𝑑 𝑒
𝑑 𝑒 𝑓 =𝑎 𝑒 𝑓
+𝑏 +𝑐
𝑕 𝑖 𝑖 𝑔 𝑔 𝑕
𝑔 𝑕 𝑖
= 𝑎 𝑒𝑖 − 𝑓𝑕 + 𝑏 𝑓𝑔 − 𝑑𝑖 + 𝑐(𝑑𝑕 − 𝑔𝑒)
The following properties hold:
(i) If two rows or two columns of a square matrix 𝑨 are identical, then 𝑑𝑒𝑡 𝑨 = 0.
(ii) If two rows or two columns of a square matrix 𝑨 are interchanged, then only the
sign of 𝑑𝑒𝑡 𝑨 is changed.
(iii) The value of 𝑑𝑒𝑡 𝑨 is unchanged if a multiple of one row is added to another row, or
if a multiple of one column is added to another column.
(iv) If 𝑨 and 𝑩 are square matrices of the same order, then 𝑑𝑒𝑡(𝑨𝑩) = (𝑑𝑒𝑡𝑨) (𝑑𝑒𝑡𝑩).

(v) If 𝑨 is invertible, then 𝑑𝑒𝑡(𝑨−𝟏 ) = (𝑑𝑒𝑡𝑨)–1 .

(vi) If 𝑨 is an 𝑛 × 𝑛 matrix, then 𝑑𝑒𝑡 𝑘𝑨 = 𝑘 𝑛 𝑑𝑒𝑡 𝑨.
DETERMINANT OF A SQUARE MATRIX: Let A = aij be a square matrix of order 𝑛.
m×n
a11 a12 ⋯ a1n
a21 a22 ⋯ a2n
Then the number ⋯ ⋯ ⋯ ⋯ is called the determinant of the matrix A and is
⋯ ⋯ ⋯ ⋯
an1 an2 ⋯ ann
denoted by A or by Det. A or by aij . Since in a determinant the number of rows is
equal to the number of columns, therefore only square matrices can have determinants.

The value of a determinant does not change when rows and columns are interchanged.

If any two rows (or two columns) of a determinant are interchanged, the value of the
determinant is multiplied by–1.

If all the elements of one rows (or one column) of a determinant are multiplied by the
same number k, the valuye of the new determinant is k times the value of the given
determinant.

If two rows (or two columns) of a determinant are identical, the value of the
determinant is zero.

In a determinant the sum of the products of the elements of any row (column) with the
cofactors of the corresponding elements of any other row (column) is zero.

If in a determinant each element in any row (or column) consists of the sum of two
terms, then the determinant can be expressed as the sum of two determinants of the
same order.

If to the elements of a row (or column) of a determinant are added m times the
corresponding elements of another row (or column), the value of the determinant thus
obtained is equal to the value of the original determinant.

DETERMINANTS OF ORDER 𝐧: A determinant of order n has n rows and n columns. It

has n × n elements.
A determinant of order n is a square array of n × n quantities (numbers of functions)
a11 a12 ⋯ a1n
a21 a22 ⋯ a2n
enclosed between vertical bars, ∆= ⋯ ⋯ ⋯ ⋯ .
⋯ ⋯ ⋯ ⋯
an1 an2 ⋯ ann

The cofactor Aij of the element aij in ∆ is equal to (−1)i+j times the determinant of order
n − 1 obtained from ∆ by leaving the row and the column passing through the element
aij . Then we have

∆= ai1 Ai1 + ai2 Ai2 + ⋯ + ain Ain i = 1,2,3, … . , or n ,

Or ∆= a1j A1j + a2j A2j + ⋯ + anj Anj j = 1,2,3, … . , or n .

DETERMINANT TRICK: Let 𝑀 be a finitely-generated module over a unital commutative

ring 𝑅, let 𝐽 be an ideal of 𝑅, and let 𝜙: 𝑀 → 𝑀 be an endomorphism of the 𝑅-module 𝑀.
Suppose that 𝜙(𝑀) ⊂ 𝐽𝑀. Then there exist elements 𝑎0 , 𝑎1 , . . . , 𝑎𝑛−1 of R such that
𝑛−1
𝑎𝑘 ∇ 𝐽 𝑛−𝑘 for 𝑘 = 1, 2, . . . , 𝑛 − 1 and 𝜙 𝑛 + 𝑘=0 𝑎𝑘 𝜙
𝑘
= 0𝐸𝑛𝑑 𝑅 (𝑀) .

DETERMINSTIC MODEL: The inventory models, in which demand is assumed to be fixed

for a subsequent period of time, are known as deterministic model.

DEVELOPABLE SURFACE: The envelope of a single parameter family of planes is called a

developable surface or simply a developable. It is a surface isometric to the plane.

DEVIATION: The difference between a value in a frequency distribution and a reference

value. If {𝑥𝑖 } is a set of observations of the random variable 𝑋 then 𝑥𝑖 – 𝑥 is the deviation
of the 𝑖th observation from the mean.

DIAGONAL FORMS: Diagonal forms are some of the simplest projective varieties to
study from an arithmetic point of view (including the Fermat varieties). Their local zeta-
functions are computed in terms of Jacobi sums. Waring's problem is the most classical
case.
DIAGONALIZABLE MATRIX: A matrix is said to be diagonalizable if it is similar to a
diagonal matrix.
An 𝑛 × 𝑛 matrix is diagonalizable if and only if it possesses 𝑛 linearly independent
eigenvectors.

DIAGONAL MATRIX: A square matrix in which all entries, except the main diagonal
entries, are zero. A square matrix A = aij those elements above and below the
n×n

principle diagonal are all zero, i.e., aij = 0 for all i ≠ j, is called a diagonal matrix

Thus a diagonal matrix is both upper and lower triangular. An n- rowed diagonal matrix
whose diagonal elements in order are d1 , d2 , d3 , … , dn will often be denoted by the
symbol

Diag [d1 , d2 , … , dn ]

For example

2 0 0 2 0 0
A = 0 0 0 and 0 2 0 are diagonal matrices.
0 0 5 0 0 2

DIAMETER OF A METRIC SPACE: Diameter of a metric space is the supremum of

distances between pairs of points.

DIAMETRAL PLANE: The locus of points which bisect a system of parallel chords of a
given sphere is called the diametric plane, all chords being parallel to a given line.

DIFFEOMORPHISM: Given two differentiable manifolds 𝑀 and 𝑁, a bijective

map 𝑓 from 𝑀 to 𝑁 is called a diffeomorphism if both 𝑓: 𝑀 → 𝑁 and its inverse 𝑓 −1 : 𝑀 →
𝑁 are smooth functions.
DIFFERENCE EQUATION: Difference equations describe the change with time of
variables that change over discrete time steps. Difference equations have some
similarities with differential equations. The difference is that the independent variable
in a differential equation can vary continuously. In a difference equation, the function
has one value at time 1, then another value at time 2, another value at time 3, and so on.

Here is an example of a difference equation:

𝑥𝑖 = 𝑎 + 𝑏𝑥𝑖−1
DIFFERENTIABLE: A continuous function is said to be differentiable over an interval
𝑎, 𝑏 if its derivative exists everywhere in that interval. This means that the graph of the
function is smooth, with no kinks, cusps, or breaks.
DIFFERENTIAL EQUATION: A differential equation is an equation containing the
derivatives of a function with respect to one or more independent variables. The order
of the equation is the highest derivative that appears.
DIFFERENTIAL GEOMETRY: The area of mathematics which uses differential calculus in
the study of geometry. For example, to prove that the area of a circle is exactly 𝜋𝑟 2 .
DIFFERENTIATION: Differentiation is the process of finding a derivative.
DIGRAPH: A digraph or directed graph consists of a number of vertices, some of which
are joined by arcs, where an arc, or directed edge, joins one vertex to another and has an
arrow on it to indicate its direction. The arc from the vertex 𝑢 to the vertex 𝑣 may be
denoted by the ordered pair (𝑢, 𝑣).
DIHEDRAL ANGLE: The angle formed by two intersecting planes. The size of the
dihedral angle is defined as the size of the angle formed by two intersecting lines (one in
each plane) that are both perpendicular to the straight edge along which the two planes
intersect.
DIHEDRAL GROUP: The group of symmetries of a regular 𝑛 −sided polygon; the notation
𝐷𝑛 is often used.
DILATION: Dilation of a map between metric spaces is the infimum of numbers L such
that the given map is L-Lipschitz.
dim: It is the abbreviation for dimension of a vector space.
DIMENSION OF A VECTOR SPACE: The number n of vectors in a basis of the finite-
dimensional vector space V is called the dimension of V and we write 𝑑𝑖𝑚(𝑉) = 𝑛. Thus,
as we might expect, ℝ𝑛 has dimension n. 𝐾 𝑥 is infinite-dimensional, but the space
𝐾 𝑥 ≤ 𝑛 of polynomials over 𝐾 of degree at most 𝑛 has basis 1, 𝑥, 𝑥 2 , … … … , 𝑥 𝑛 so its
dimension is (𝑛 + 1).
Note that the dimension of 𝑉 depends on the field 𝐾. Thus the complex numbers ℂ
can be considered as

 A vector space of dimension 1 over ℂ ,with one possible basis being the single
element 1.
  A vector space of dimension 2 over ℝ, with one possible basis given by the two
elements 1,i
 A vector space of infinite dimension over ℚ.
Note the following:
 Suppose that the vectors 𝑣1 , 𝑣2 , ⋯ ⋯ , 𝑣𝑛 , 𝑤 span 𝑉 and that 𝑤 is a linear
combination of 𝑣1 , 𝑣2 , ⋯ ⋯ , 𝑣𝑛 .Then 𝑣1 , 𝑣2 , ⋯ ⋯ , 𝑣𝑛 span 𝑉.
 Suppose that the vectors 𝑣1 , 𝑣2 , ⋯ ⋯ , 𝑣𝑟 span the vector space 𝑉.Then there is a
subsequence of 𝑣1 , 𝑣2 , ⋯ ⋯ , 𝑣𝑟 which forms a basis of 𝑉.
 Let 𝑉 be the vector space over 𝐾 which has a finite spanning set, and suppose
that the vectors 𝑣1 , 𝑣2 , ⋯ ⋯ , 𝑣𝑟 are linearly independent in 𝑉. Then we can extend
the sequence to a basis 𝑣1 , 𝑣2 , ⋯ ⋯ , 𝑣𝑛 of 𝑉 where 𝑛 ≥ 𝑟.
 Suppose that the vectors v𝟏 , v𝟐 , ⋯ ⋯ , v𝒏 span 𝑉 and that vectors
w𝟏 , w𝟐 , ⋯ ⋯ , w𝒎 ϵ V are linearly independent. Then 𝑚 ≤ 𝑛.
 Let V be a vector space of dimension n over K. Then any n vectors which span V
forms a basis of V, and no n-1 vectors can span V.
 Let V be a vector space of dimension n over K. Then any linearly independent
vectors form a basis of V and no 𝑛 + 1 vectors can be linearly independent.
 If a non-trivial vector space V is spanned by a finite number of vectors, then it
has a basis.

DIMENSION THEOREM FOR VECTOR SPACES: Given a vector space 𝑉, any two linearly
independent generating sets (in other words, any two bases) have the same cardinality.
If 𝑉 is finitely generated, then it has a finite basis, and the result says that any two bases
have the same number of elements.

DINI'S THEOREM (ANALYSIS): Dini's theorem says that if a monotone sequence of

functions converges on a compact space, it converges uniformly. If 𝑋 is
a compact topological space, and { 𝑓𝑛 } is a monotonically increasing
sequence of continuous real-valued functions on 𝑋 which converges pointwise to a
continuous function 𝑓, then the convergence is uniform. The same conclusion holds if
{ 𝑓𝑛 } is monotonically decreasing instead of increasing.

DINI’S THEOREM (DIFFERENTIAL GEOMETRY): Two surface which are mapped

geodescially on each other by a non-conformal mapping must have line elements which
can be written in the forms
𝑑𝑠 2 = 𝑈 − 𝑉 (𝑑𝑢2 + 𝑑𝑣 2 )

𝑑𝑠 ∗2 = 𝑉 −1 − 𝑈 −1 (𝑈 −1 + 𝑑𝑢2 + 𝑉 −1 𝑑𝑣 2 )

Where 𝑈 = 𝑈 𝑢 𝑎𝑛𝑑 𝑉 = 𝑉 𝑣 .

DIOPHANTINE EQUATION: An algebraic equation in one or more unknowns, with

integer coefficients, for which integer solutions are required. A great variety of
Diophantine equations have been studied. Some have infinitely many solutions, some
have finitely many and some have no solutions. For example:
(i) 14𝑥 + 9𝑦 = 1 has solutions 𝑥 = 2 + 9𝑡, 𝑦 = – 3– 14𝑡 .
(ii) 𝑥 2 + 1 = 2𝑦 4 has two solutions 𝑥 = 1, 𝑦 = 1 and 𝑥 = 239, 𝑦 = 13.
(iii) 𝑥 4 + 𝑦 3 = 𝑧 5 has no solutions.
DIOPHANTUS OF ALEXANDRIA: Diophantus of Alexandria was an ancient Greek
mathematician whose work displayed an algebraic approach to the solution of
equations in one or more unknowns, unlike earlier Greek methods that were more
geometrical. In the books of Arithmetica that survive, particular numerical examples of
more than 100 problems are solved, probably to indicate the general methods of
solution. These are mostly of the kind now referred to as Diophantine equations.
DIRAC DELTA FUNCTION: We have functions which have non-zero values on very short
intervals. The Dirac delta function may be thought of as a generation of this concept. The
Dirac delta function and its derivatives play a very useful role in the solution of the
boundary value problems in mathematical physics as well as in quantum mechanics.

Consider the function

1
, t >𝜖
δϵ t = 2ϵ ………….. 1
0, t <𝜖

∞ 𝜖 1
Thus ∫−∞ δϵ t dt = ∫−𝜖 2ϵ dt. ………….. 2

Again if 𝑓 𝑡 is any function which is integrable in the interval −ϵ, ϵ then using the
mean value theorem of the integral calculus, we have

∞ 1 𝜖
∫−∞ 𝑓(𝑡)δϵ t dt = 2ϵ ∫−𝜖 𝑓(𝑡)dt = 𝑓 𝜃𝜖 , 𝜃 ≤ 1. ……….…. 3

Thus we may think of a limiting function denoted by δ t approached by δϵ t as ϵ → 0.

i.e., δ t = limϵ→0 δϵ(t) …………… 4

as ϵ → 0, from (1) and (2), we have

∞ if t = 0
δ t = limϵ→0 δϵ(t) = ………….. 5
0 if t ≠ 0
∞
and ∫−∞ δ t dt = 1. ………… 6

this limiting function δ t defined by equations (5) and (6) is known a the Dirac delta
function (or the unit impulse function) after Dirac who first introduced it. Dirac called
this delta function as improper function as there cannot be proper function satisfying
such conditions.

DIRECTED GRAPH: An directed graph or digraph (𝑉, 𝐸) consists of a finite set 𝑉

together with a subset 𝐸 of 𝑉 × 𝑉 . The elements of 𝑉 are the vertices of the digraph;
the elements of 𝐸 are the edges of the digraph.
DIRECTED NUMBER: A number having a positive or negative sign attached, such as +4
or -4. It is commonly used for temperature readings to emphasize whether it is above or
below 0°.
DIRECTIONAL DERIVATIVE: The directional derivative of a function 𝑓 (𝑥, 𝑦) in the
direction of a unit vector 𝒗 = (𝑣𝑥 , 𝑣𝑦 ) is the dot product of the gradient of 𝑓 with 𝒗:
𝜕𝑓 𝜕𝑓
𝑑𝑖𝑟𝑒𝑐𝑡𝑖𝑜𝑛𝑎𝑙 𝑑𝑒𝑟𝑖𝑣𝑎𝑡𝑖𝑣𝑒 = 𝑣𝑥 + 𝑣
𝜕𝑥 𝜕𝑦 𝑦
DIRECTRIX: A directrix is a line that helps to define a geometric figure.
DIRECT SUM OF MODULES: Let 𝑀1 , 𝑀2 , . . . , 𝑀𝑘 be modules over a unital commutative
ring 𝑅. The direct sum 𝑀1 ⊕ 𝑀2 ⊕ · · · ⊕ 𝑀𝑘 is defined to be the set of ordered 𝑘-tuples
(𝑥1 , 𝑥2 , . . . , 𝑥𝑘 ), where 𝑥𝑖 ∇ 𝑀𝑖 for 𝑖 = 1, 2, . . . , 𝑘. This direct sum is itself an 𝑅-module:
(𝑥1 , 𝑥2 , . . . , 𝑥𝑘 ) + (𝑦1 , 𝑦2 , . . . , 𝑦𝑘 ) = (𝑥1 + 𝑦1 , 𝑥2 + 𝑦2 , . . . , 𝑥𝑘 + 𝑦𝑘 ), 𝑟(𝑥1 , 𝑥2 , . . . , 𝑥𝑘 ) =
(𝑟𝑥1 , 𝑟𝑥2 , . . . , 𝑟𝑥𝑘 ) for all 𝑥𝑖 , 𝑦𝑖 ∇ 𝑀𝑖 and 𝑟 ∇ 𝑅. If 𝐾 is any field, then 𝐾 𝑛 is the direct
sum of 𝑛 copies of 𝐾.
DIRECT SUM OF SUBSPACES: If 𝑈 ∩ 𝑊 = {0}, then we write 𝑈 ⊕ 𝑊 for 𝑈 + 𝑊. A
vector space 𝑉 is said to be the direct sum of subspaces 𝑈 and 𝑊 if 𝑉 = 𝑈 ⊕ 𝑊.
DIRICHLET, PETER GUSTAV LEJEUNE (1805–59): Peter Gustav Lejeune Dirichlet was a
German mathematician who was professor at the University of Berlin before succeeding
Gauss at the University of Göttingen. He proved that in any arithmetic series
𝑎, 𝑎 + 𝑑, 𝑎 + 2𝑑, …, where 𝑎 and 𝑑 are relatively prime, there are infinitely many
primes. He gave the modern definition of a function. In more advanced work, he was
concerned to see analysis applied to number theory and mathematical physics.
DIRICHLET’S PRIME NUMBER THEOREM: Suppose that 𝑞 ∇ 𝑁 and 𝑎 ∇ 𝑍 satisfy
(𝑎, 𝑞) = 1. Then there are infinitely many primes 𝑝 ≡ 𝑎 𝑚𝑜𝑑 𝑞.
DIRICHLET SERIES: A series in the form
∞

𝑎𝑛 𝑒 −𝜆 𝑛 𝑧
𝑛 =1

where 𝑎𝑛 and 𝑧 are complex and {𝜆𝑛 } is a monotonic increasing sequence of real
numbers. When 𝜆𝑛 = log 𝑒 𝑛, the series reduces to
∞

𝑎𝑛 𝑛−𝑧
𝑛=1

known as the Dirichlet L-series.

DIRICHLET’S CONDITIONS: Let 𝑓 𝑥 be a function which satisfies the following
conditions:

(a) 𝑓 𝑥 is defined on some interval say 𝑥𝜖 −𝜆, 𝜆 .

(b) 𝑓 𝑥 𝑎𝑛𝑑 𝑓′ 𝑥 are piecewise continuous in −𝜆, 𝜆 .
(c) 𝑓 𝑥 is periodic whose period is 2𝜆.

The above mentioned conditions are known as Dirichlet’s Conditions.

DIRICHLET’S TEST: A test for convergence of a series. If 𝑎𝑛 is a series which has

bounded partial sums,
i.e.
𝑚

𝑎𝑛 < 𝐾
𝑛=1

for all values of 𝑚, and {𝑏𝑛 } is decreasing sequence converging to zero, then 𝑎𝑛 𝑏𝑛
converges.
DISCONNECTED GRAPH: D is connected graph is a graph in which the vertices are
separated into two or more distinct groups and cannot be linked from a vertex in one
group to a vertex in the other group through a series of edges.
DISCRETE: Discrete refers to a situation where the possibilities are distinct and
separated from each other. For example, the number of people in a city is discrete,
because there is no such thing as a fractional person.
DISCRETE RANDOM VARIABLE: A discrete random variable is a random variable which
can only take on values from a discrete list. The probability function (or density
function) lists the probability that the variable will take on each of the possible values.
The sum of the probabilities for all of the possible values must be 1. Binomial
distribution, Poisson distribution, Geometric distribution, Hypergeometric distribution
etc are examples of discrete random variable.
DISCRETE TOPOLOGY: Let 𝑋 be any non-empty set and let 𝑃(𝑋) be the collection of all
subsets of 𝑋. Then 𝑃(𝑋) is called the discrete topology on the set 𝑋. The topological
space (𝑋; 𝑃(𝑋)) is called a discrete space.
DISCRETE VALUATION: Let 𝐾 be a field, and let 𝑍 ∪ {∞} be the set obtained from the
ring 𝑍 of integers by adding a symbol ∞ with the properties that ∞ + ∞ = ∞, 𝑛 +
∞ = ∞ + 𝑛 = ∞, ∞ − 𝑛 = ∞ 𝑎𝑛𝑑 ∞ > 𝑛 for all integers 𝑛. A discrete valuation on
the field 𝐾 is a function 𝜈: 𝐾 → 𝑍 ∪ {∞} which satisfies the following conditions:
(i) 𝜈(𝑎) = ∞ if and only if 𝑎 = 0𝐾 ;
(ii) 𝜈(𝑎𝑏) = 𝜈(𝑎) + 𝜈(𝑏) for all 𝑎, 𝑏 ∇ 𝐾;
(iii) 𝜈(𝑎 + 𝑏) ≥ 𝑚𝑖𝑛(𝜈(𝑎), 𝜈(𝑏)) for all 𝑎, 𝑏 ∇ 𝐾.
Let 𝑝 be a prime number. Then, given any non-zero rational number 𝑟, there exist
integers 𝑘, 𝑢 and v such that 𝑟 = 𝑝𝑘 𝑢𝑣 −1 and neither 𝑢 nor 𝑣 is divisible by 𝑝. The
integer 𝑘 is uniquely determined by 𝑟, and we define 𝑣𝑝 (𝑟) = 𝑘. We also define
𝑣𝑝 (0) = ∞. Then the function 𝑣𝑝 ∶ 𝑄 → 𝑍 ∪ {∞} defined in this fashion is a discrete
valuation on the field 𝑄 of rational numbers.
DISCRETE VALUATION RING: A discrete valuation ring is an integral domain 𝑅 with a
unique maximal ideal 𝑀 whose proper ideals are all of the form 𝑀𝑘 for some positive
integer 𝑘. Every discrete valuation ring is a principal ideal domain.
DISCRIMINANT: The discriminant D of a quadratic equation 𝑎𝑥 2 + 𝑏𝑥 + 𝑐 = 0 is
𝐷 = 𝑏 2 + 4𝑎𝑐. If 𝑎, 𝑏 and 𝑐 are real numbers, the discriminant allows us to determine
the characteristics of the solution for 𝑥. If 𝐷 is a positive perfect square, then 𝑥 will have
two rational values. If 𝐷 = 0, then 𝑥 will have two real and equal solutions. If 𝐷 is
positive but is not a perfect square, then 𝑥 will have two irrational solutions. If 𝐷 is
negative, then 𝑥 will have two complex solutions. In other words, the equation has two
distinct real roots, equal roots (that is, one root) or no real roots according to whether
the discriminant is positive, zero or negative.
DISJOINT SETS: Two sets are said to be disjoint if they have no elements in common,
that is, if their intersection is the empty set.
DISJOINT UNION: The disjoint union 𝑈 ⊔ 𝑉 is the union of two disjoint sets , 𝑉 ; the
notation 𝑈 ⊔ 𝑉 simply means 𝑈 ∪ 𝑉 together with an assertion (or reminder) that 𝑈
and 𝑉 are disjoint. The disjoint union of any collection of (disjoint) sets is defined and
denoted similarly.
DISJUNCTION: A disjunction is an OR statement of the form: “A OR B.” It is true if either
A or B is true.
DISPERSION: The spread of values of a variable in a distribution. A measure of
dispersion is a way of describing how scattered or spreads out the observations in a
sample are. The term is also applied similarly to a random variable. Common measures
of dispersion are the range, interquartile range, mean absolute deviation, variance and
standard deviation. The range may be unduly affected by odd high and low values. The
mean absolute deviation is difficult to work with because of the absolute value signs.
The standard deviation is in the same units as the data, and it is this that is most often
used. The interquartile range may be appropriate when the median is used as the
measure of location.
DISPLACEMENT VECTOR: A vector that describes a movement from one point to
another, using both direction and magnitude.
DISTANCE BETWEEN TWO CODEWORDS: The distance between two codewords in a
binary code is the number of bits in which the two codewords differ. For example, the
distance between 010110 and 001100 is 3 because they differ in the second, third and
fifth bits. If the distance between any two different codewords in a binary code is at
least 3, the code is an error-correcting code capable of correcting any one error.
DISTANCE FROM A POINT TO A LINE IN THE PLANE: If a point 𝑃 has coordinates (𝑥1 , 𝑦1 )
and a line 𝑙 has equation 𝑎𝑥 + 𝑏𝑦 + 𝑐 = 0, then the distance from 𝑃 to 𝑙 is equal to
𝑎𝑥1 + 𝑏𝑦1 + 𝑐
𝑎2 + 𝑏 2
DISTANCE FROM A POINT TO A PLANE IN 3-DIMENSIONAL SPACE: If a point 𝑃 has
coordinates (𝑥1 , 𝑦1 , 𝑧1 ) and a plane 𝑝 has equation 𝑎𝑥 + 𝑏𝑦 + 𝑐𝑧 + 𝑑 = 0, the
distance from 𝑃 to 𝑝 is equal to
 𝑎𝑥1 + 𝑏𝑦1 + 𝑐𝑧1 + 𝑑
𝑎2 + 𝑏 2 + 𝑐 2
DISTANCE OF THE CODE: Let 𝐶 be a code. The distance of the code, denoted 𝑑(𝐶), is
defined by 𝑑(𝐶) = 𝑚𝑖𝑛{𝑑(𝑐1 , 𝑐2 ) | 𝑐1 , 𝑐2 ∇ 𝐶, 𝑐1 , ≠ 𝑐2 }. An (𝑛, 𝑀) −code of distance 𝑑 is
called an (𝑛, 𝑀, 𝑑) −code. The values 𝑛, 𝑀, 𝑑 are called the parameters of the code.
DISTRIBUTIVE PROPERTY: The distributive property says that
𝑎 𝑏 + 𝑐 = 𝑎𝑏 + 𝑎𝑐 ∀ 𝑎, 𝑏 𝑎𝑛𝑑 𝑐.
DIVERGENCE: The divergence of a vector field 𝒇 (written as ∆. 𝒇) is defined to be the
scalar
𝜕𝑓𝑥 𝜕𝑓𝑦 𝜕𝑓𝑧
∆. 𝒇 = + +
𝜕𝑥 𝜕𝑦 𝜕𝑧
It can be thought of as the dot product of the operator ∆ (del) with the field f.
DIVERGENCE THEOREM: The divergence theorem states that if E is a three-dimensional
vector field, then the surface integral of E over a closed surface is equal to the triple
integral of the divergence of E over the volume enclosed by that surface:

𝐸. 𝑑𝑆 = (∆. 𝐸)𝑑𝑉

DIVERGENT SERIES: A divergent series is an infinite series with no finite sum.

DIVISION ALGEBRA: A division algebra is an algebra with identity which, as a ring, is a
division ring.

DIVISION ALGORITHM: For integers 𝑎 and 𝑏, with 𝑏 > 0, there exist unique integers 𝑞
and 𝑟 such that 𝑎 = 𝑏𝑞 + 𝑟, where 0 ≤ 𝑟 < 𝑏. In the division of 𝑎 by 𝑏, the number 𝑞
is the quotient and 𝑟 is the remainder.
DIVISOR OF ZERO: If in a ring there are non-zero elements 𝑎 and 𝑏 such that 𝑎𝑏 = 0,
then 𝑎 and 𝑏 are divisors of zero. For example, in the ring of 2 × 2 real matrices,
0 1 0 0 0 0
=
0 0 1 0 0 0
and so each of the matrices on the left-hand side is a divisor of zero. In the ring 𝑍6 ,
consisting of the set 0, 1, 2, 3, 4, 5 with addition and multiplication modulo 6, the
element 4 is a divisor of zero since 4.3 = 0.
DNE: It is the abbreviation for a solution for an expression does not exist, or is
undefined. Generally used with limits and integrals.
DOMAIN: The domain of a function is the set of all possible values for which the function
is well defined.
DOMAIN OF A RELATION: The set of all first coordinates of the ordered pairs.
DOMAIN OF A VARIABLE: The base set of numbers or quantities that can be mapped to a
second set. In elementary algebra, the domain of a function 𝑦 = 𝑓(𝑥) is the set of
values that the independent variable 𝑥 can take. For example, if 𝑦 = 𝑓(𝑥) = 𝑎𝑟𝑐𝑠𝑖𝑛(𝑥),
then the domain might be defined as all the numbers whose absolute value is no more
than 1, that is, |𝑥| < 1.
DOMAIN WITH SMOOTH BOUNDARY: Let 𝑀 be an 𝑚-dimensional abstract manifold. A
nonempty open subset 𝐷 ⊂ 𝑀 is said to be a domain with smooth boundary if for each
𝑝 ∇ 𝜕𝐷 there exists an open neighborhood 𝑊 of 𝑝 in 𝑀, and a smooth function
𝑓: 𝑊 → 𝑅 such that 0 is a regular value, 𝑓(𝑝) = 0 and
𝑊 ∩ 𝐷 = {𝑥 ∇ 𝑊 | 𝑓(𝑥) > 0}.
DORN’S CONVERSE DUALITY THEOREM: Let 𝐶 be positive semi- definite. If x, u solves
Quadratic Duality Problem, then some x ∇ Rn (not necessarily equal to x), satisfying
𝐶 x − x = 0 solves Quadratic maximization Problem and the two extremes are equal.

DOT PRODUCT: Let a and b be two n-dimensional vectors, whose components are:

𝒃 = ( 𝑏1 , 𝑏2 , 𝑏3 , − − − 𝑏𝑛 )
𝒂 = ( 𝑎1 , 𝑎2 , 𝑎3 , − − − 𝑎𝑛 )
Then
𝒃. 𝒂 = 𝑏1 𝑎1 + 𝑏2 𝑎2 + 𝑏3 𝑎3 + − − − + 𝑏𝑛 𝑎𝑛
DOUBLE-ANGLE FORMULA: The trigonometric functions of a double-angle that are
expressed in terms of separate functions of the angle.

𝑠𝑖𝑛(2𝑥) = 2 𝑠𝑖𝑛𝑥 𝑐𝑜𝑠𝑥

𝑐𝑜𝑠(2𝑥) = 𝑐𝑜𝑠2𝑥 − 𝑠𝑖𝑛2𝑥
𝑡𝑎𝑛(2𝑥) = 2𝑡𝑎𝑛𝑥/(1 − 𝑡𝑎𝑛2𝑥)

DOUBLE FAMILY OF CURVES: The quadratic differential equation of the form

𝑃𝑑𝑢² + 2𝑄𝑑𝑢𝑑𝑣 + 𝑅𝑑𝑢𝑑𝑣 + 𝑅𝑑𝑣² = 0

Where P,Q, R are continuous functions of 𝑢 𝑎𝑛𝑑 𝑣 and do not vanish together represents
two families of curves on the surface provided 𝑄 2 − 𝑃𝑅 > 0.

We can write equation (1) in the form

 2
𝑑𝑢 𝑑𝑢
𝑃 + 2𝑄 +𝑅 =0
𝑑𝑣 𝑑𝑣

𝑑𝑢
Which is quadratic in 𝑑𝑣 and by solving this equation the separate differential equations

for the two families are obtained.

DOUBLE INTEGRAL: The double integral of a two-variable function 𝑓 (𝑥, 𝑦) represents

the volume under the surface 𝑧 = 𝑓 (𝑥, 𝑦) and above the 𝑥𝑦 plane in a specified region.
For example:
𝑥=𝑏,𝑦=𝑑

𝑓 𝑥, 𝑦 𝑑𝑥𝑑𝑦
𝑥=𝑎,𝑦=𝑐

represents the volume under the surface 𝑧 = 𝑓 (𝑥, 𝑦) over the rectangle from 𝑥 = 𝑎 to
𝑥 = 𝑏 and 𝑦 = 𝑐 to 𝑦 = 𝑑.
DOUBLE-NAPPED CONE: Two identical but opposite cones that share a common vertex.
DOUBLE SERIES: A sequence with the indices, i.e., a mapping from the Cartesian
product 𝑁 × 𝑁 of two copies of the set of natural numbers 𝑁 to a subset of the real or
complex numbers, is called a double sequence and is denoted by 𝑎𝑚𝑛 or 𝑎𝑚 ,𝑛 . If
there exists a number 𝐼 such that for any positive 𝜀, there is a natural number 𝑁(𝜀)
satisfying 𝑎𝑚𝑛 − 1 < 𝜀 for all 𝑚 > 𝑁(𝜀) and 𝑛 > 𝑁(𝜀), then we say that the sequence
𝑎𝑚𝑛 has a limit 𝐼 and write lim𝑚 →∞ 𝑎𝑚𝑛 = 𝛼𝑛 uniformly in 𝑛 and lim𝑚 →∞,𝑛→∞ 𝑎𝑚𝑛 =
∞
𝐼. For a given double sequence 𝑎𝑚𝑛 , the formal series m,n=1 𝑎𝑚𝑛 is called a double
series and is sometimes denoted by 𝑎𝑚𝑛 . In contrast with double series, the ordinary
series discussed previously is called a simple series.

DOUBLING TIME: The time it takes for a population to double itself.

DUAL BASIS: Let 𝑉 be a finite-dimensional real vector space, let 𝑢1 , 𝑢2 , . . . , 𝑢𝑛 be a basis of
𝑉 . The corresponding dual basis of the dual space 𝑉 ∗ of 𝑉 consists of the linear
𝑛
functionals 𝜀1 , 𝜀2 , . . . , 𝜀𝑛 on 𝑉 , where ϵi 𝑗 =1 𝜆𝑗 𝑢𝑗 = 𝜆𝑖 for 𝑖 = 1, 2, . . . , 𝑛 and for all
real numbers 𝜆1 , 𝜆2 , . . . , 𝜆𝑛 .
DUAL CODE: Let 𝐶 be a linear code over 𝐹𝑞𝑛 . Then
1. The dual code of 𝐶 is 𝐶 ⊥ (the orthogonal complement of 𝐶 in 𝐹𝑞𝑛 ).
2. The dimension of 𝐶 is the dimension of 𝐶 as a vector subspace of 𝐹𝑞𝑛 , and is denoted
dim(C).
DUALITY (FUNCTIONAL ANALYSIS): If 𝑋 and 𝑌 are normed linear spaces and
𝑇 ∶ 𝑋 → 𝑌 , then we get a natural map 𝑇 ∗ : 𝑌 ∗ → 𝑋 ∗ by 𝑇 ∗ 𝑓(𝑥) = 𝑓(𝑇 𝑥) for all
𝑓 ∇ 𝑌 ∗ , 𝑥 ∇ 𝑋. In particular, if 𝑇 ∇ 𝐵(𝑋, 𝑌 ), then 𝑇 ∗ ∇ 𝐵(𝑌 ∗ , 𝑋 ∗ ). In fact,
𝑇 ∗ 𝐵(𝑌 ∗ , 𝑋 ∗ ) = 𝑇 𝐵(𝑋, 𝑌 ). To prove this, note that |𝑇 ∗ 𝑓(𝑥)| = |𝑓(𝑇 𝑥)| ≤
𝑓 𝑇 𝑥 . Therefore 𝑇 ∗ 𝑓 ≤ 𝑓 𝑇 , so 𝑇 ∗ is indeed bounded, with 𝑇 ∗ ≤ 𝑇 .
Also, given any 𝑦 ∇ 𝑌 , we can find 𝑔 ∇ 𝑌 ∗ such that |𝑔(𝑦)| = 𝑦 , 𝑔 = 1. Applying
this with 𝑦 = 𝑇 𝑥 (𝑥 ∇ 𝑋 arbitrary), gives 𝑇 𝑘 = |𝑔(𝑇 𝑥)| = |𝑇 ∗ 𝑔𝑥| ≤
𝑇 ∗ 𝑘𝑔𝑘 𝑥 = 𝑇∗ 𝑥 . This shows that 𝑇 ≤ 𝑇 ∗ . Note that if 𝑇 ∇ 𝐵(𝑋, 𝑌 ), 𝑈 ∇
𝐵(𝑌, 𝑍), then (𝑈𝑇)∗ = 𝑇 ∗ 𝑈 ∗.

DUALITY IN LINEAR PROGRAMMING: Every linear programming problem is associated

with another linear programming problem called the dual of the problem. The original
problem is called ‘primal’ while the other is called its dual. The solution of the dual
problem leads to the solution of the primal problem and thus efficient computational
techniques can be developed through the concept of duality.

Primal Problem: Find 𝑥𝑗 ≥ 0 (𝑗 = 12, … . , 𝑛) in order to determine

𝑧 = 𝑐1 𝑥2 + 𝑐2 𝑥2 + ⋯ + 𝑐𝑛 𝑥𝑛 ≤ 𝑏1

Subject to the constraints

𝑎11 𝑥1 + 𝑎12 𝑥2 + ⋯ + 𝑎1𝑛 𝑥𝑛 ≤ 𝑏1

𝑎21 𝑥1 + 𝑎22 𝑥2 + ⋯ + 𝑎2𝑛 𝑥𝑛 ≤ 𝑏2

⋮ ⋮ ⋮

𝑎𝑚1 𝑥1 + 𝑎𝑚2 𝑥2 + ⋯ + 𝑎𝑚𝑛 ≤ 𝑏𝑚

The corresponding dual problem is obtained by transposing the rows and columns of
constraint coefficients, transporting the coefficients of the objectives function and the
right hand side of the constraints, reversing the inequalities and minimizing instead of
maximizing.

DUALITY SPECULATION: Every primal linear programming problem has associated with
it another linear programming problem called the dual linear programming problem.
 Constraint parameters in one are variable coefficients in other.
  Coefficients for the objective function of either are the RHS for other problem.
DUALITY THEOREM:
 Feasible solutions exist and objective function is bounded, then same is true for
other problem.
 Feasible solutions exist and objective function is unbounded, then other problem
is infeasible.
 No feasible solutions exist, then other problem is either infeasible or has
unbounded objective function.
Some properties of this theory are:
Weak duality If 𝒙 is a feasible solution for the primal problem and 𝒚 is a
property feasible solution for the dual problem, then 𝒄𝒙 ≤ 𝒚𝒃.
Strong duality If 𝒙∗ is an optimal solution for the primal problem and 𝒚∗ is an
property optimal solution for the dual, then 𝒄𝒙∗ = 𝒚∗ 𝒃.
Complementary At each iteration, the simplex method simultaneously identifies
solutions property a CPF solution 𝒙 for the primal problem and a complementary
solution 𝒚 for the dual (found in row 0, the coefficient of the
slack variables), where 𝒄𝒙 = 𝒚𝒃.
Complementary At the final iteration, the simplex method simultaneously
optimal solutions identifies an optimal solution 𝒙∗ for the primal problem and a
property complementary optimal solution 𝒚∗ for the dual problem (found
in row 0, the coefficient of the slack variables), where
𝒄𝒙∗ = 𝒚∗ 𝒃.. In this solution, 𝒚∗ gives the shadow prices for the
primal problem.
Symmetry property The dual of the dual is the primal.
Complementary basic Each basic solution in the primal problem has a complementary
solutions property basic solution in the dual, where their respective objective
function values (Z and W) are equal.
Complementary The variables in the primal basic solution and the
slackness property complementary dual basic solution satisfy the complementary
slackness as shown:
Primal variable Dual variable
Basic Non-basic (m variables)
 Non-basic Basic (n variables)

Complementary Each optimal basic solution in the primal problem has a

optimal basic complementary optimal basic solution in the dual problem,
solutions property where their respective objective function values (Z and W ) are
equal.

Relationships between Primal Basic solution and Dual Basic Solution:

Primal Basic Soln Feasible Dual Basic Soln Feasible

suboptimal Yes superoptimal no
optimal Yes optimal yes
superoptimal No suboptimal yes
infeasible, not sup No infeasible, not sup no

Shortcut for conversion between primal and dual:

Primal Dual
Maximize Z Minimize W
Constraint i: Variable yi:
≤ form yi ≥ 0
= form unconstrained
≥ form y’i ≤ 0
Variable xj: Constraint j:
xj ≥ 0 ≥ form
unconstrained unconstrained
x’j ≤ 0 ≤ form
DUAL OF A QUOTIENT SPACE: Next, consider the projection map 𝜋 ∶ 𝑋 → 𝑋/𝑆 where 𝑆
is a closed subspace. We then have 𝜋 ∗ ∶ (𝑋/𝑆)∗ → 𝑋 ∗ . Since 𝜋 is surjective, this map is
injective. It is easy to see that the range is contained in 𝑆 𝑎 . In fact we now show that 𝜋 ∗
maps (X/S) ∗ onto 𝑆 𝑎 , hence provides a canonical isomorphism of 𝑆 𝑎 with (𝑋/𝑆)∗ .
Indeed, if 𝑓 ∇ 𝑆 𝑎 , then we have a splitting 𝑓 = 𝑔 ∘ 𝜋 with 𝑔 ∇ (𝑋/𝑆)∗ (just define
𝑔(𝑐) = 𝑓(𝑥) where x is any element of the coset c). Thus 𝑓 = 𝜋 ∗ 𝑔 is indeed in the
range of 𝜋 ∗ . This correspondence is again an isometry.

DUAL OF A SUBSPACE: An important case is when 𝑇 is the inclusion map 𝑖 ∶ 𝑆 → 𝑋,

where 𝑆 is a closed subspace of 𝑋. Then 𝑟 = 𝑖 ∗ ∶ 𝑋 ∗ → 𝑆 ∗ is just the restriction map:
𝑟𝑓 (𝑠) = 𝑓(𝑠). Hahn-Banach tells us that 𝑟 is surjective. Obviously 𝑁 (𝑟) = 𝑆 𝑎 . Thus
we have a canonical isomorphism 𝑟 ∶ 𝑋 ∗ /𝑆 𝑎 → 𝑆 ∗ . In fact, the Hahn-Banach theorem
shows that it is an isometry. Via this isometry one often identifies 𝑋 ∗ /𝑆𝑎 with 𝑆 ∗ .

DUAL PROBLEM: Find 𝑦𝑗 ≥ 0 (𝑗 = 1,2, … , 𝑚) in order to minimize

𝑧 = 𝑏1 𝑦1 + 𝑏2 𝑦2 + … + 𝑏𝑚 𝑦𝑚

Subject to the constraints

𝑎11 𝑦1 + 𝑎21 𝑦2 + … . . +𝑎𝑚1 𝑦𝑚 ≥ 𝑐1

𝑎12 𝑦1 + 𝑎22 𝑦2 + ⋯ + 𝑎𝑚2 𝑦𝑚 ≥ 𝑐2
⋮ ⋮ ⋮ ⋮
𝑎1𝑛 𝑦1 + 𝑎2𝑛 𝑦2 + ⋯ + 𝑎𝑛𝑚 𝑦𝑚 ≥ 𝑐𝑚

DUAL SIMPLEX METHOD: For a L.P.P. (maximization problem) the optimality criterion
of the simplex method 𝑐𝑗 − 𝑧𝑗 = 𝑐𝑗 − 𝑐𝐵 𝐵 −1 . 𝛼𝑗 ≤ 0, for all 𝑗 where 𝐵 is the basis,
depends only on 𝛼𝑗 , and 𝑐𝑗 and is independent of the requirement vector 𝑏. Thus, every
basic solution with all 𝑐𝑗 − 𝑧𝑗 ≤ 0, will not be feasible but any basic feasible solution
with all 𝑐𝑗 − 𝑧𝑗 ≤ 0 will certainly be an optimal solution.

DUAL SPACE (FUNCTIONAL ANALYSIS): The dual space X* of a normed space X is the set
𝑠𝑢𝑝
of continuous linear functionals on X. Define a norm on it by 𝛼 = 𝑥 ≤ 1 𝛼(𝑥)

DUAL SPACE (LINEAR ALGEBRA): For fixed vector spaces 𝑈 and 𝑉 over 𝐾, the
operations of addition and scalar multiplication on the set 𝐻𝑜𝑚𝐾 (𝑈, 𝑉) of all linear
maps from 𝑈 to 𝑉 makes 𝐻𝑜𝑚𝐾 (𝑈, 𝑉) into a vector space over 𝐾.

Given a vector space 𝑈 over a field 𝐾, the vector space 𝑈 ∗ = 𝐻𝑜𝑚𝐾 (𝑈, 𝑉) plays a
special role.It is often called the dual space or the space of convectors of 𝑈. One can
think of coordinates as elements of 𝑈 ∗ .Indeed,let 𝑒𝑖 be a basis of U. Every x 𝜖 𝑈 can be
uniquely written as x =𝛼1 𝑒1 + 𝛼2 𝑒2 + ⋯ ⋯ + 𝛼𝑛 𝑒𝑛 , 𝛼𝑖 𝜖 𝐾
The elements 𝛼𝑖 depends on 𝑥 as well as on a choice of the basis, so for each I one can
write the coordinate function 𝑒 𝑖 : 𝑈 → 𝐾 , 𝑒 𝑖 𝑥 = 𝛼𝑖

It is easy to check that 𝑒 𝑖 is a linear map, and indeed the function 𝑒 𝑖 form a basis of the
dual space 𝑈 ∗ .

DUAL SPHERICAL TRIANGLES: Let ∆ be a spherical triangle with angles 𝛼, 𝛽, 𝛾 and side
lengths 𝑎, 𝑏, 𝑐. Then the dual triangle ∆∗ has sides of length 𝑎∗ = 𝜋 − 𝛼, 𝑏 ∗ = 𝜋 − 𝛽, 𝑐 ∗ =
𝜋 − 𝛾 and angles 𝛼 ∗ = 𝜋 − 𝑎, 𝛽 ∗ = 𝜋 − 𝑏, 𝛾 ∗ = 𝜋 − 𝑐.
DYADIC TENSOR: A dyadic tensor has order two, and may be represented as a square
matrix. The conventions aij, aij, and aij, do have different meanings (the position of the
index determines its valence (variance), in that the first may represent a quadratic
form, the second a linear transformation, and the differention is significant in contexts
that require tensors that aren't orthogonal. A dyad is a tensor such as aibj, product
component-by-component of rank-one tensors. In this case it represents a linear
transformation, of rank one in the sense of linear algebra - a clashing terminology that
can cause confusion.
DYNAMIC PROGRAMMING: The area of mathematics relating to the study of
optimization problems where a step-wise decision making approach is employed. This
is often done iteratively.

E
ECCENTRICITY: The ratio of the distances between a point on a conic and a fixed point
(the focus) and between the point and a fixed line (the directrix). For 𝑒 = 0 a circle is
produced, for 0 < 𝑒 < 1 an ellipse is produced, for 𝑒 = 1 a parabola is produced and
for 𝑒 > 1 the conic produced is a hyperbola.
ECHELON FORM OF A MATRIX: A matrix is in echelon form if
 All the zero rows come below the non-zero rows.
 The first non-zero entry in each non-zero row is 1 and occurs in a column to the
right of the leading 1 in the row above.
EDGE: The edge of a polyhedron is a line segment where two faces intersect.
𝑎. 𝑒
EGOROV THEOREM: If 𝑓𝑛 𝑓 on a finite measure set 𝑋 then for any 𝜍 > 0 there is
→
𝐸𝜍 ⊂ 𝑋 with µ(𝐸𝜍 ) < 𝜍 and 𝑓𝑛 ⇒ 𝑓 on 𝑋 ∖ 𝐸𝜍 .
𝜇 𝑎. 𝑒
If 𝑓𝑛 𝑓 then there is a subsequence (𝑛𝑘 ) such that 𝑓𝑛 𝑘 𝑓 for 𝑘 → ∞.
→ →

Indicator function: For 𝐴 ⊆ 𝑋, we define 𝜒𝐴 to be the indicator function of 𝐴, by

1: 𝑥 ∇ 𝐴
𝜒𝐴 𝑥 = . Then, if 𝜒𝐴 is measurable, then 𝜒𝐴 −1 ( (1/2,3/2) ) = 𝐴 ∇ 𝐿;
0: 𝑥 ∈ 𝐴
conversely, if 𝐴 ∇ 𝐿, then 𝑋 ∖ 𝐴 ∇ 𝐿, and we see that for any 𝑈 ⊆ ℝ open, 𝜒𝐴 −1 (𝑈) is
either ∅, 𝐴, 𝑋 ∖ 𝐴, or 𝑋, all of which are in 𝐿. So 𝜒𝐴 is measurable if and only if 𝐴 ∇ 𝐿.

Ei: It is the abbreviation for exponential integral function.

EIGENVALUE: Suppose that a square matrix 𝑨 multiplies a vector 𝒙, and the resulting
vector is proportional to 𝒙:
𝐴𝑥 = 𝜆𝑥
In this case, l is said to be an eigenvalue of the matrix A, and x is the corresponding
eigenvector. In order to find the eigenvalues, rewrite the equation like this:
𝐴 − 𝜆𝐼 𝑥 = 𝑂
where 𝑰 is the appropriate identity matrix.
EIGENVALUE OF OPERATOR: An eigenvalue of operator 𝑇 ∇ 𝐵(𝐻) is a complex number
𝜆 such that there exists a nonzero 𝑥 ∇ 𝐻, called eigenvector with property 𝑇𝑥 = 𝜆 𝑥, in
other words 𝑥 ∇ 𝑘𝑒𝑟(𝑇 − 𝜆 𝐼).

 In finite dimensions 𝑇 − 𝜆 𝐼 is invertible if and only if 𝜆 is not an eigenvalue.

 In infinite dimensions it is not the same: the right shift operator S is not
invertible but 0 is not its eigenvalue because 𝑆𝑥 = 0 implies 𝑥 = 0.
EINSTEIN, ALBERT (1879–1955): Albert Einstein was an outstanding mathematician.
Born in Ulm in Germany, he lived in Switzerland and Germany before moving to the
United States in 1933. He was responsible in 1905 for the Special Theory of Relativity
and in 1916 for the General Theory. He made a fundamental contribution to the birth of
quantum theory and had an important influence on thermodynamics. He regarded
himself as a physicist rather than as a mathematician, but his work has triggered off
many developments in modern mathematics.
EISENSTEIN’S IRREDUCIBILITY CRITERION: Let 𝑓(𝑥) = 𝑎0 + 𝑎1 𝑥 + 𝑎2 𝑥 2 + · · · + 𝑎𝑛 𝑥 𝑛
be a polynomial of degree 𝑛 with integer coefficients, and let 𝑝 be a prime number.
Suppose that
• 𝑝 does not divide 𝑎𝑛 ,
• 𝑝 divides 𝑎0 , 𝑎1 , . . . , 𝑎𝑛−1 ,
• 𝑝2 does not divide 𝑎0 .
Then the polynomial 𝑓 is irreducible over the field 𝑄 of rational numbers.
EINSTEIN NOTATION: This notation is based on the understanding that in a product of
two indexed arrays, if an index letter in the first is repeated in the second, then the
interpretation is that the product is summed over all values of the index. For example if
aij is a matrix, then under this convention aii is its trace. The Einstein convention is
generally used in physics and engineering texts, to the extent that if summation is not to
be applied it is normal to note that explicitly.
EINSTEIN’S SUMMATION CONVENTION: Einstein proposed that if the same index letter
appears twice in a term, then it will automatically be assumed to be summed over. For
example, 𝑎𝑖 𝑥 𝑖 means the summation 𝑎𝑖 𝑥 𝑖 .
ELEMENTARY ABELIAN GROUP: An Abelian group all non-trivial elements of which have
the same prime order p.
ELEMENTARY AXIOM SYSTEM: An axiom system written in the language of restricted
predicate calculus. Examples of elementary axiom systems are those of formal
arithmetic, the Zermelo–Fraenkel system of set theory and the system of the theory of
types.
ELEMENTARY COLUMN OPERATION: Elementary column operations are the following
operations on the columns of a matrix:
 Interchange two columns,
 Multiply a column by a non-zero scalar,
  Add a multiple of one column to another column.
An elementary column operation can be produced by post-multiplication by the
appropriate elementary matrix.
ELEMENTARY DIVISORS OF A MATRIX OVER A POLYNOMIAL RING: Powers of the monic
irreducible polynomials over the field 𝐾 into which the invariant factors of 𝐹[𝑥] split.
Two 𝑚 × 𝑛-matrices over 𝐾[𝑥] having the same rank are equivalent (that is, can be
obtained from one another by means of elementary operations) if and only if they have
the same system of elementary divisors.
ELEMENTARY FUNCTION: following real functions are called elementary functions:
The rational functions, the trigonometric functions, the logarithmic and exponential
functions, the functions 𝑓 defined by 𝑓(𝑥) = 𝑥 𝑚 /𝑛 (where m and n are non-zero
integers), and all those functions that can be obtained from these by using addition,
subtraction, multiplication, division, composition and the taking of inverse functions.
ELEMENTARY INTERVAL OF A PARTIALLY ORDERED SET: A subset consisting of two
elements 𝒂 ≤ 𝒃 such that there are no other elements in the partially ordered
set between them, i.e.
𝑎 ≤ 𝑥 ≤ 𝑏 ⇒ 𝑎 = 𝑥 𝑜𝑟 𝑎 = 𝑏.

ELEMENTARY MATRIX: Elementary matrix is a square matrix obtained from the identity
matrix 𝐼 by an elementary row operation. A matrix obtained from a unit matrix by a
single elementary transformation is called an elementary matrix (or 𝐸- matrix). For
0 0 1 1 0 0 1 2 0
example 0 1 0 , 0 4 0 , 0 1 0 are the elementary matrices obtained from
1 0 0 0 0 1 0 0 1
Is by subjecting it to the elementary operations C1 ⟷ C3 , R 2 → 4R 2 , R1 → R1 + R 2
respectively.

The elementary matrix corresponding to the 𝐸-operation R i ⟷ R j is its own inverse.

The inverse of the 𝐸- matrix corresponding to the 𝐸-operation R1 → kR i , (k ≠ 0), is

the 𝐸- matrix corresponding to the 𝐸-operation R i → k −1 R1 .

The inverse of the 𝐸- matrix corresponding to the 𝐸-operation R i → R i + kR j is the 𝐸-

matrix corresponding to the 𝐸-operation R1 → R i − kR j .
ELEMENTARY NUMBER THEORY: The branch of number theory that investigates
properties of the integers by elementary methods. These methods include the use of
divisibility properties, various forms of the axiom of induction and combinatorial
arguments. Sometimes the notion of elementary methods is extended by bringing in the
simplest elements of mathematical analysis. Traditionally, proofs are deemed to be non-
elementary if they involve complex numbers. Usually, one refers to elementary number
theory the problems that arise in branches of number theory such as the theory of
divisibility, of congruences, of arithmetic functions, of indefinite equations, of partitions,
of additive representations, of the approximation by rational numbers, and of continued
fractions. Quite often, the solution of such problems leads to the need to go beyond the
framework of elementary methods. Occasionally, following the discovery of a non-
elementary solution of some problem, one also finds an elementary solution of it.
ELEMENTARY OPERATION: Addition, subtraction, multiplication, division and finding
integer roots are the elementary operations.
ELEMENTARY ROW OPERATION: Much type of calculations with matrices can be carried
out in a computationally efficient manner by the use of certain type of operations on
rows and columns. These are really the same as the operations used in solving sets of
simultaneous linear equations.

Let 𝐴 be an 𝑚 × 𝑛 matrix over 𝐾 with rows 𝑟1 , 𝑟2 ,⋯ ⋯ ⋯ , 𝑟𝑚 𝜖 𝐾 1,𝑛 . The three types of

elementary row operations on A are defined as follows:

(R1) for some i≠ 𝑗 ,add a multiple of 𝑅𝑗 𝑡𝑜 𝑅𝑖

3 1 9 𝑅3 →𝑅3 −3𝑅1 3 1 9
Example: 4 6 7 4 6 7
2 5 8 −7 2 −19

(R2) Interchange two rows.

(R3) Multiply a row by a non-zero scalar

2 0 5 𝑅2 →2𝑅2 2 0 5
Example: 1 −2 3 2 −4 6
5 1 2 5 1 2

An elementary row operation can be produced by post-multiplication by the

appropriate elementary matrix.
ELEMENTARY TRANSFORMATIONS OF A MATRIX: An elementary transformation (or an
𝐸- Transformation) is an operation of any one of the following types:

1. The interchange of any two rows (or columns).

2. The multiplication of the elements of any row (or column) be any non-zero
number.
3. The addition to the elements of any other row (or column) the corresponding
elements of any other row (or column) multiplied by any number.

An elementary transformation is called a row transformation or a column

transformation according as it applies to rows or columns.

ELEMENT OF BEST APPROXIMATION: An element 𝒖𝟎 in a given set 𝑭 that is a best

approximation to a given element 𝒙 in a metric space 𝑿, i.e. is such that
𝜌(𝒖𝟎 , 𝑥) = 𝑖𝑛𝑓{𝜌(𝑢, 𝑥): 𝑥 ∇ 𝐹} .

This is a generalization of the classical concept of a polynomial of best approximation.

The main questions concerning elements of best approximation are: their existence and
uniqueness, their characteristic properties, the properties of the operator that
associates with each element 𝑥 ∇ 𝑋 the set of elements of best approximation and
numerical methods for the construction of elements of best approximation.

ELLIPSE: An ellipse is the set of all points in a plane such that the sum of the distances
to two fixed points is a constant. Each of these two fixed points is known as a focus or
focal point. The plural of focus is foci. The longest distance across the ellipse is known as
the major axis. The shortest distance across is the minor axis. The center of the ellipse
is the midpoint of the segment that joins the two foci. The equation of an ellipse with
center at the origin is
𝑥2 𝑦2
+ =1
𝑎2 𝑏 2
where 2𝑎 is the length of the major axis, and 2𝑏 is the length of the minor axis. Area of
this ellipse is 𝜋𝑎𝑏.
ELLIPSOID: A quadric whose equation in a suitable coordinate system is
𝑥2 𝑦2 𝑧2
+ + =1
𝑎2 𝑏 2 𝑐 2
The three axial planes are planes of symmetry. All non-empty plane sections are
ellipses. If the ellipse has major axis 2𝑎 and minor axis 2𝑏, then the ellipsoid formed by
4
rotating the ellipse about its major axis will have the volume 𝜋𝑎𝑏 2 .
3

ELLIPSOIDAL COORDINATES: The numbers 𝜆, 𝜇 and 𝜈 connected with Cartesian

rectangular coordinates 𝑥, 𝑦 and 𝑧 by the formulas

2
𝜆 + 𝑎2 𝜇 + 𝑎2 𝜈 + 𝑎2
𝑥 =
𝑏 2 − 𝑎2 𝑐 2 − 𝑎2

2
𝜆 + 𝑏2 𝜇 + 𝑏2 𝜈 + 𝑏2
𝑦 =
𝑎2 − 𝑏 2 𝑐 2 − 𝑏 2

𝜆 + 𝑐2 𝜇 + 𝑐2 𝜈 + 𝑐2
𝑧2 =
𝑎2 − 𝑐 2 𝑏 2 − 𝑐 2

where −𝑎2 < 𝜈 < −𝑏 2 < 𝜇 < −𝑐 2 < 𝜆 < ∞ .

ELLIPTIC COORDINATES: Two numbers 𝜍 and 𝜏 connected with rectangular Cartesian

coordinates by the formulas

𝜍 + 𝑎2 𝜏 + 𝑎2
𝑥2 =
𝑎2 − 𝑏 2

𝜍 + 𝑏2 𝜏 + 𝑏2
𝑦2 =
𝑏 2 − 𝑎2
where −𝑎2 < 𝜏 < −𝑏 2 < 𝜍 < ∞ .

ELLIPTIC CYLINDER: Elliptic cylinder is a cylinder in which the fixed curve is an ellipse
and the fixed line to which the generators are parallel is perpendicular to the plane of
the ellipse. It is a quadric, and in a suitable coordinate system has equation
ELLIPTIC FUNCTION: A function defined on the complex plane for which 𝑓(𝑧) = 𝑓(𝑧 +
𝑎) = 𝑓(𝑧 + 𝑏) where 𝑎/𝑏 is not real. From this it follows that 𝑓(𝑧 + 𝑚𝑎 + 𝑛𝑏) =
𝑓(𝑧) for all integers 𝑚, 𝑛 and that the function is periodic in two distinct directions on
the complex plane.

A doubly-periodic analytic function 𝑓 𝑧 is said to be an elliptic function if its only

possible singular points in the finite part of the plane are poles.

All the meromorphic functions which have two distinct periods are called elliptic
functions.

Primitive period parallelogram, Mess and cell: Let 𝑓 𝑧 be an elliptic function with 2ω1
and 2ω2 as a pair of primitive periods.

ω2
Supposing that imaginary part of which is not zero to be positive, we see that the
ω1

points 0, 2ω1 , 2ω1 + 2ω2 , 2ω1 taken in order are the vertices of a parallelogram
described in the positive sense. This is called the primitive period-parallelogram of the
elliptic function 𝑓 𝑧 . There are the unlimited number of such primitive period
parallelograms. Now the only periods of 𝑓 𝑧 (double periodic function) are of the form
2mω1 + 2nω2 , where m and n are integers. Therefore the vertices are the only points
within or on a primitive period- parallelogram. Whose affixes are periods in the Argand
plane the points of affix 2mω1 + 2nω2 where m = 0, ±1, ±2, … , n = 0, ±1, ±2…. Are
denoted by Ωm,n

i.e. Ωm,n = 2mω1 + 2nω2

Thus, the four points Ωm,n, Ωm+1,n Ωm+1,n+1 are the vertices of a parallelogram. This
parallelogram may be obtained from the primitive period-parallelogram by a
translation without rotation. This parallelogram with vertices
Ωm′n′ Ωm+1′n Ωm+1,n+1′ Ωm ′ n+1 is called a period parallelogram or a mesh. The Argand
plane may be covered by this system of non overlapping meshes. In each mesh (within
and on the boundary) there are only a finite number of poles and zeros, and hence we
can translate the mesh without rotation until no pole or zero lies on its boundary. The
parallelogram thus obtained is called a cell.

The set of poles (or zeros) in a given cell is called an irreducible set.

ELLIPTIC GEOMETRY: A geometry in a space with a Riemannian curvature that is

constant and positive in any two-dimensional direction. Elliptic geometry is a higher-
dimensional generalization of the Riemann geometry.
ELLIPTIC INTEGRALS: The integrals

𝜔 𝑑𝑡
𝑧 = ∫0 𝑘 <1 ……. i
1−𝑡 2 1−𝑘 2 𝑡 2

is called an elliptic integral of the first kind. The integral exists if ω is real and such that
ω < 1. By analytical continuation it can be extended to other values of ω.

If 𝑡 = sin 𝜃 𝑎𝑛𝑑 ω = sin ∅, the integral (i) assumes an equivalent form

∅ 𝑑𝑡
𝑧 = ∫0 ………. ii
1−𝑘 2 𝑠𝑖𝑛 2 𝜃

Where we often write ∅ = 𝑎𝑚 𝑧.

Now if 𝑘 = 0, from (i), we have

𝜔
𝑑𝑡
= 𝑠𝑖𝑛−1 𝜔
0 1− 𝑡2

∴ 𝜔 = 𝑠𝑖𝑛 𝑍.

∴ when 𝑘 ≠ 0, the integral (i) is deputed by Sn −1 (𝜔/𝑘) or briefly Sn −1 (𝜔), when k doe
not change during a given discussion.

Thus we have

𝜔 𝑑𝑡
𝑧 = Sn −1 𝜔 = ∫0 1−𝑡 2 1−𝑘 2 𝑡 2

This gives the function Sn (z), which is called a Jacobi’s elliptic function.

ELLIPTIC PARABOLOID: A quadric whose equation in a suitable coordinate system is

 𝑥 2 𝑦 2 2𝑧
+ =
𝑎2 𝑏 2 𝑐
Here the 𝑦𝑧-plane and the 𝑧𝑥-plane are planes of symmetry. Sections by planes 𝑧 = 𝑘,
where 𝑘 ≥ 0, are ellipses (circles if 𝑎 = 𝑏); planes 𝑧 = 𝑘, where 𝑘 < 0, have no
points of intersection with the paraboloid. Sections by planes parallel to the 𝑦𝑧-plane
and to the 𝑧𝑥-plane are parabolas. Planes through the 𝑧-axis cut the paraboloid in
parabolas with vertex at the origin.
ELLIPTIC POINTS (DIFFERENTIAL GEOMETRY): The points on the surface at which the
principal curvatures 𝜅𝑎 𝑎𝑛𝑑 𝜅𝑏 have the same sign. i.e. the Gaussian curvature 𝜅 is
positive are called elliptic points.

𝐿𝑁−𝑀² 𝐿𝑁−𝑀²
Again 𝐾 = 𝜅𝑎 𝜅𝑏 = =
𝐸𝐺−𝐹² 𝐻²

Thus, a point is an elliptic point if 𝐿𝑁 − 𝑀² > 0

EMBEDDED PARAMETRIZED MANIFOLD: A regular parametrized manifold 𝜍: 𝑈 → 𝑅 𝑛

which is a homeomorphism 𝑈 → 𝜍(𝑈), is called an embedded parametrized manifold.
In particular this definition applies to curves and surfaces, and thus we can speak of
embedded parametrized curves and embedded parametrized surfaces. In addition to
being smooth and regular, the condition on 𝜍 is thus that it is injective and that the
inverse map 𝜍(𝑥) → 𝑥 is continuous 𝜍(𝑈) → 𝑈.
EMPIRICAL STATEMENT: A statement that is based upon observation and experimental
evidence.
EMPTY SET: Empty set also called void set or null set is the set, denoted by 𝜑 with no
elements in it. Consequently, its cardinality, 𝑛(𝜑), is zero. The Zermelo-Fraenkel axioms
of set theory postulate that there exists an empty set.
ENCRYPT: To encrypt means to transform information or data into a coded form.
END: It is the abbreviation for categories of endomorphisms.

ENDOMORPHISM RING: The associative ring 𝑬𝒏𝒅 𝑨 = 𝑯𝒐𝒎 (𝑨, 𝑨) consisting of all
morphisms of into itself, where is an object in some additive category. The
multiplication in 𝑬𝒏𝒅 𝑨 is composition of morphisms, and addition is the addition of
morphisms defined by the axioms of the additive category. The identity morphism 𝟏𝑨 is
the unit element of the ring 𝑬𝒏𝒅 𝑨. An element 𝝋 in 𝑬𝒏𝒅 𝑨 is invertible if and only
if 𝝋 is an automorphism of the object 𝑨. If 𝑨 and 𝑩 are objects of an additive category 𝑪,
then the group 𝑯𝒐𝒎 (𝑨, 𝑩) has the natural structure of a right module over 𝑬𝒏𝒅 𝑨 and
of a left module over 𝑬𝒏𝒅 𝑩. Let 𝑻: 𝑪 → 𝑪𝟏 be a covariant (or contravariant) additive
functor from an additive category 𝑪 into an additive category 𝑪𝟏 . Then for any
object 𝑨 in 𝑪 the functor 𝑻 induces a natural homomorphism (or anti-
homomorphism) 𝑬𝒏𝒅 𝑨 → 𝑬𝒏𝒅 𝑻(𝑨).
END-POINT: End point of an interval is a number defining one end of an interval on the
real line. Each of the finite intervals [𝑎, 𝑏], (𝑎, 𝑏), [𝑎, 𝑏) and (𝑎, 𝑏] has two end-points, 𝑎
and 𝑏. Each of the infinite intervals [𝑎, ∞), (𝑎, ∞), (– ∞, 𝑎] and (– ∞, 𝑎) has one end-
point, 𝑎.
ENERGY EQUATION: The principle of energy enunciates that the change in the total
energy (Kinetic 𝑇, Potential 𝑊 and Intrinsic 𝐼) is equal to work done by the extraneous
forces. The potential due to external forces is supposed to the independent of time.

∂q 1
The Euler’s equation of motion is = −∆Ω − ρ ∆ρ, where the extraneous forces are
∂t

conservative.

∂q
ρ = −ρ ∆Ω − ∆ρ
∂t
Multiplying the above equation scalarly by q, we get

∂q 1 d
ρq ∂t = −ρq ∆Ω − q∆ρ or ρ dt q2 + ρq. ∆Ω = −q. ∆ρ
2

1 d dΩ d 1
or ρ dt q2 + ρ dt = −q. ∆ρ or ρ dt q2 + Ω − q. ∆ρ .
2 2

Integrating the relation (2) over 𝑉, we have

d 1 2
ρ q + Ω dv = − q. ∆ρ dv
V dt 2 V

Or
𝑑 1 2
ρq dv + ρΩ dv = − q. ∆ρ dv
𝑑𝑡 V2 V V

𝑑
Or 𝑇 + 𝑊 = − ∫V ∆. ρq dv + ∫V ρ ∆. q dv
𝑑𝑡

𝑑 ρ dρ
𝑇 + 𝑊 = − ∫V ∆. ρq dv − ∫V dv.
𝑑𝑡 ρ dt
By virtue of divergence theorem and the equation of continuity,

The R.H.S. may be expressed as

𝑑 dI
𝑇+𝑊 =− ρq ∙ n ds −
𝑑𝑡 s dt

𝑑
Or 𝑇 + 𝑊 + 𝐼 †= − ∫s ρq ∙ n ds,
𝑑𝑡

which is known as energy equation.

ENERGY PRINCIPLE: Denote by 𝐸 the class of all measures of finite energy. A symmetric
kernel is called positive definite (or of positive type) if 𝜇 − 𝑣, 𝜇 − 𝑣 = 𝜇, 𝜇 + 𝑣, 𝑣 −
2(𝜇, 𝑣) ≥ 0 for any 𝜇, 𝑣𝜖 𝐸. If the equality (𝜇 − 𝑣, 𝜇 − 𝑣) = 0 always implies 𝜇 = 𝑣, then
the kernel is said to satisfy the energy principle.

ENTIRE FUNCTION: A function which has no singularity in the finite part of plane is
called an entire function.

ENUMERATION PROBLEM: An algorithmic problem in which one has to construct an

algorithm that enumerates 𝑨 for a given set 𝑨, i.e. an algorithm 𝕳 that is applicable to
any natural number, that converts it to an element of 𝑨 and such that any element
of 𝑨 is obtained by applying 𝕳 to some natural number; in other words, 𝑨 =
𝕳 𝒊 : 𝒊 ∇ 𝚴 . The enumeration problem for a set 𝑨 is solvable if and only if 𝑨 is a non-
empty enumerable set.
ENVELOPE: A curve or surface that is tangential to every curve or surface in a family.
For example, if a family of circles has radius 𝑎 and centre at a distance 𝑟 > 𝑎 from a
fixed point 𝐶, the envelope will be an annulus with the two circles having radii 𝑟 + 𝑎
and 𝑟 – 𝑎.
ENVELOPING SERIES FOR A NUMBER A: A series
∞

𝑎𝑛
𝑛=𝑜

such that

|𝐴 − (𝑎0 + ⋯ + 𝑎𝑛 )| < |𝑎𝑛+1 |

for all 𝑛 = 0,1, …. An enveloping series may converge or diverge; if it converges, then its
sum is equal to 𝐴.
EPICYCLOID: The curve traced out by a point on the circumference of a circle rolling
round the outside of a fixed circle. When the two circles have the same radius, the curve
is a cardioid. A planar curve given by the trajectory of a point on a circle rolling on the
exterior side of another circle. The parametric equations are:
𝜃
𝑥 = (𝑟 + 𝑅)𝑐𝑜𝑠𝜃 − 𝑟𝑐𝑜𝑠[(𝑟 + 𝑅) ],
𝑟

𝜃
𝑦 = (𝑟 + 𝑅)𝑠𝑖𝑛𝜃 − 𝑟𝑠𝑖𝑛[(𝑟 + 𝑅) ],
𝑟

where r is the radius of the rolling and 𝑅 that of the fixed circle, and 𝜃 is the angle
between the radius vector of the point of contact of the circles and the x-axis.

EPIMORPHISM: Epimorphism is a type of morphism 𝑓: 𝑋 → 𝑌 with the property that

for all morphisms g between 𝑌 and 𝑍, 𝑔1 ∘ 𝑓 = 𝑔2 ∘ 𝑓 ⇒ 𝑔1 = 𝑔2 .
EPSILON: The Greek letter 𝑒, written ∇, commonly used to represent a small, strictly
positive quantity.
EPSILON–DELTA NOTATION: The standard notation used to define the concepts of
limits and continuity of functions, sequences, series, nets, filters etc.
EQUALITY (FUNCTIONS): Functions 𝑓 and 𝑔 are said to be equal if they have the same
domains and for every 𝑥 in their common domain, 𝑓 𝑥 = 𝑔 𝑥 .
EQUALITY (FUZZY SETS): Let 𝐴 and 𝐵 are fuzzy subsets of a classical set 𝑋. 𝐴 and 𝐵 are
said to be equal, denoted 𝐴 = 𝐵, if 𝐴 ⊂ 𝐵 and 𝐵 ⊂ 𝐴 in the sense of Fuzzy
subsethood. Note that 𝐴 = 𝐵 if and only if 𝜇𝐴 (𝑥) = 𝜇𝐵 (𝑥) ∀ 𝑥 ∇ 𝑋.
EQUALITY (MATRICES): Matrices 𝑨 and 𝑩, where 𝑨 = [𝑎𝑖𝑗 ] and 𝑩 = [𝑏𝑖𝑗 ], are said to
be equal if and only if they have the same order and 𝑎𝑖𝑗 = 𝑏𝑖𝑗 for all 𝑖 and 𝑗.
EQUALITY (SETS): Sets 𝐴 and 𝐵 are said to be equal if they consist of the same elements.
In order to establish that 𝐴 = 𝐵, a technique that can be useful is to show, instead, that
both 𝐴 ⊆ 𝐵 and 𝐵 ⊆ 𝐴.
EQUALITY OF POLYNOMIALS: Two matric polynomials are equal iff (if and only if), the
coefficients of the power of λ are the same. Every square whose elements are ordinary
polynomials in λ, can essentially be expressed as a matrix polynomial in λ of degree m is
the highest power of λ occurring in any element of the matrix.

EQUAL VECTORS: Two vectors are said to be equal if, and only if, they are parallel, have
the same sense of direction and the same magnitude. The starting points of the vectors
are immaterial. It is the direction and magnitude which are important. To denote the
equality of vectors, the usual equality sign (=) is used. Thus, if 𝒂 and 𝒃 are equal
vectors, we write 𝒂 = 𝒃.

EQUATING COEFFICIENTS: Let 𝑓(𝑥) and 𝑔(𝑥) be polynomials, and let

𝑓(𝑥) = 𝑎0 + 𝑎1 𝑥 + … + 𝑎𝑛−1 𝑥 𝑛 −1 + 𝑎𝑛 𝑥 𝑛 ,
𝑔(𝑥) = 𝑏0 + 𝑏1 𝑥 + … + 𝑏𝑛−1 𝑥 𝑛−1 + 𝑏𝑛 𝑥 𝑛 ,
where it is not necessarily assumed that 𝑎𝑛 ≠ 0 and 𝑏𝑛 ≠ 0. If 𝑓(𝑥) = 𝑔(𝑥) for all
values of 𝑥, then 𝑎0 = 𝑏0 , 𝑎1 = 𝑏1 , … , 𝑎𝑛−1 = 𝑎𝑛−1 , 𝑎𝑛 = 𝑎𝑛 . Using this fact is known
as equating coefficients.
EQUATING REAL AND IMAGINARY PARTS: Complex numbers 𝑎 + 𝑏𝑖 and 𝑐 + 𝑑𝑖 are
equal if and only if 𝑎 = 𝑐 and 𝑏 = 𝑑. Using this fact is called equating real and
imaginary parts. For example, if (𝑎 + 𝑏𝑖)2 = 3 + 22𝑖, then 𝑎2 – 𝑏 2 = 3 and
2𝑎𝑏 = 22.
EQUATION: An equation is a statement that asserts that two mathematical expressions
are equal in value. We can also say that an equation is a statement that says that two
mathematical expressions have the same value. If this is true for all values of the
variables involved then it is called an identity, for example 2(𝑥 – 3) = 2𝑥 – 6, and
where it is only true for some values it is called a conditional equation; for example
𝑥 2 – 2𝑥 – 3 = 0 is only true when 𝑥 = – 1 𝑜𝑟 3, which are known as the roots of the
equation. An equation that can be put in the general form 𝑎𝑥 + 𝑏 = 0, where 𝑥 is
unknown and 𝑎 and 𝑏 are known, is called a linear equation. Any one-unknown equation
can be written in this form provided that it contains no terms with 𝑥 2 , 1/𝑥, or any term
with 𝑥 raised any power other than 1. An equation involving 𝑥 2 and 𝑥 is called a
quadratic equation, and can be written in the form 𝑎𝑥 2 + 𝑏𝑥 + 𝑐 = 0.
EQUATION OF A CIRCLE: The equation of a circle in the Argand plane can be put in the
form

zz + bz + bz + 𝑐 = 0

where 𝑐 is real and 𝑏 is complex constant.

EQUIANGULAR SPIRAL: Equiangular spiral is a curve whose equation in polar

coordinates is 𝑟 = 𝑎𝑒 𝑘𝜃 , where 𝑎 (> 0), 𝑘 are constants. Let 𝑂 be the origin and 𝑃 be
any point on the curve. The curve derives its name from the property that the angle 𝛼
between 𝑂𝑃 and the tangent at 𝑃 is constant. In fact, 𝑘 = 𝑐𝑜𝑡 𝛼. The equation can be
written 𝑟 = 𝑘𝜃 + 𝑏, and the curve is also called the logarithmic spiral.

EQUICONTINUITY OF A SET OF FUNCTIONS: An idea closely connected with the concept

of compactness of a set of continuous functions. Let 𝑋 and 𝑌 be compact metric spaces
and let 𝐶(𝑋, 𝑌) be the set of continuous mappings of 𝑋 into 𝑌. A set 𝐷 ⊂ 𝐶(𝑋, 𝑌) is called
equicontinuous if for any 𝜖 > 0 there is a 𝛿 > 0 such that 𝜌𝑋 (𝑥1 , 𝑥2 ) ≤ 𝛿 implies
𝜌𝑌 𝑓 𝑥1 , 𝑓 𝑥2 ≤ 𝜖 ∀ 𝑥1 , 𝑥2 ∇ 𝑋, 𝑓 ∇ 𝐷. Equicontinuity of 𝐷 is equivalent to the
relative compactness of 𝐷 in 𝐶(𝑋, 𝑌), equipped with the metric
𝜌(𝑓, 𝑔) = max 𝜌𝑌 (𝑓(𝑥), 𝑔(𝑥)) ;
𝑥∇𝑋

this is the content of the Arzelà–Ascoli theorem. The idea of equicontinuity can be
transferred to uniform spaces.

EQUICONVERGENT SERIES: Convergent or divergent series 𝑎𝑛 and 𝑏𝑛 whose

difference is a convergent series with zero sum: 𝑎𝑛 − 𝑏𝑛 . If their difference is only a
convergent series, then the series are called equiconvergent in the wide sense.
EQUIVALENCE CLASS: For an equivalence relation ~ on a set S, an equivalence class [𝑎]
is the set of elements of 𝑆 equivalent to 𝑎; that is to say, [𝑎] = {𝑥 |𝑥 ∇ 𝑆 and a ~ x}. It
can be shown that if two equivalence classes have an element in common, then the two
classes are, as sets, equal. The collection of distinct equivalence classes having the
property that every element of 𝑆 belongs to exactly one of them is a partition of 𝑆.
EQUIVALENCE OF CODES: Two (𝑛, 𝑀)-codes are equivalent if one can be derived from
the other by a permutation of the coordinates and multiplication of any specific
coordinate by a non-zero scalar. Note that permuting the coordinates or multiplying by
a non-zero scalar makes no difference to the parameters of the code. In this sense, the
codes are therefore equivalent. Every linear code C is equivalent to a linear code C′ with
a generator matrix in standard form.
EQUIVALENCE OF REAL QUADRATIC FORMS: Two real quadratic forms X T AX and Y T BY
are said to be real equivalent, orthogonally equivalent, or complex equivalent according
as there exists a non-singular real, orthogonal, or a non-singular complex matrix P such
that 𝐵 = 𝑃−1 𝐴𝑃.

EQUIVALENCE RELATION: A binary relation ~ on a set 𝑆 that is reflexive, symmetric and

transitive. For an equivalence relation ~, 𝑎 is said to be equivalent to 𝑏 when 𝑎 ~ 𝑏. It is
an important fact that, from an equivalence relation on 𝑆, equivalence classes can be
defined to obtain a partition of 𝑆. equivalence relation. The equivalence class of an
element 𝑎 ∇ 𝐴 is the set of all elements 𝑏 such that 𝑎 ∼ 𝑏. An equivalence relation on
𝐴 partitions 𝐴 into a disjoint union of equivalence classes.
EQUIVALENT MATRICES A AND B OVER A RING R: Matrices such that 𝑨 can be
transformed into 𝑩 by a sequence of elementary row-and-column transformations, that
is, transformations of the following three types:
a) permutation of the rows (or columns);
b) addition to one row (or column) of another row (or column) multiplied by an
element of R
c) multiplication of a row (or column) by an invertible element of 𝑹. Equivalently, 𝑩 is
obtained from 𝑨 by multiplication on left or right by a sequence of matrices each of
which is either a) a permutation matrix; b) an elementary matrix; c) an
invertible diagonal matrix.
EQUIVALENT NORMS: Let 𝑁 be a normed linear space. Suppose that ∥ 1 and ∥ 2 are
two different norms 𝑁. The norms ∥ 1 is said to be equivalent to ∥ 2 written as
∥ 1 ~∥ 2, if there exist positive number 𝑎 and 𝑏 such that

𝑎 𝑓 1 ≤ 𝑓 2 ≤ 𝑏 𝑓 1 , for all 𝑓 𝜖 𝑁.

ERATOSTHENES OF CYRENE (275–195 BC): Eratosthenes was a Greek astronomer and

mathematician who was the first to calculate the size of the Earth by making
measurements of the angle of the Sun at two different places a known distance apart.
His other achievements include measuring the tilt of the Earth’s axis. He is also credited
with the method known as the sieve of Eratosthenes.
ERDOS AND TURAN CONJECTURE: Every set 𝐴 of natural numbers with positive upper
asymptotic density contains arbitrarily long arithmetic progrssions. This is equivalent
to the statement that if there is a natural number 𝑡 such that the set 𝐴 contains no
arithmetic progression of 𝑡 terms, then 𝑑(𝐴) = 0.
ERDOS, PAUL (1913–96): Paul Erdos was a Hungarian mathematician who was prolific
in his output, publishing more than 1500 papers, many of them jointly as he
collaborated with a wide range of other mathematicians. He sought elegant and simple
solutions to complex problems, which requires insight into the essential nature of the
problem, as well as the technical mathematics to extract the solution. He can be
regarded as the founder of discrete mathematics and was awarded the Wolf Prize in
1983 for contributions to number theory, combinatorics, probability, set theory and
mathematical analysis.
erf: It is the abbreviation for error function.
erfc: It is the abbreviation for complementary error function.
ERROR-CORRECTING AND ERROR-DETECTING CODE: A code is said to be error-
detecting if any one error in a codeword results in a word that is not a codeword, so that
the receiver knows that an error has occurred. A code is error-correcting if, when any
one error occurs in a codeword, it is possible to decide which codeword was intended.
Certain error-correcting codes may not be able to detect errors if more than one error
occurs in a codeword; other error-correcting codes can be constructed that can detect
and correct more than one error in a codeword.
ESSENTIAL SINGULARITIES: These singularities are due to the nature (geometrical
features) of the surface and are independent of the choice of parametric representation
of the surface, e.g., the vertex of the cone is an essential singularity.
ESTIMATOR: A statistic used to estimate the value of a population parameter. An
estimator 𝑋 of a parameter 𝜃 is consistent if the probability of the difference between
the two exceeding an arbitrarily small fixed number tends to zero as the sample size
increases indefinitely. An estimator 𝑋 is an unbiased estimator of the parameter 𝜃 if
𝐸(𝑋) = 𝜃, and it is a biased estimator if not. The best unbiased estimator is the
unbiased estimator with the minimum variance. The relative efficiency of two unbiased
estimators 𝑋 and 𝑌 is the ratio 𝑉𝑎𝑟(𝑌)/𝑉𝑎𝑟(𝑋) of their variances. Estimators may be
found in different ways, including the method of maximum likelihood, and the method
of moments.
EUCLID (ABOUT 300 BC): Euclid was an outstanding mathematician of Alexandria,
author of what may well be the second most influential book in Western Culture: the
Elements. Little is known about Euclid himself, and it is not clear to what extent the
book describes original work and to what extent it is a textbook. The Elements develops
a large section of elementary geometry by rigorous logic starting from ‘undeniable’
axioms. It includes his proof that there are infinitely many primes, the Euclidean
algorithm, the derivation of the five Platonic solids, and much more.
EUCLIDEAN ALGORITHM: It is a process, based on the Division Algorithm, for finding the
greatest common divisor (𝑎, 𝑏) of two positive integers 𝑎 and 𝑏. Assuming that 𝑎 > 𝑏,
write 𝑎 = 𝑏𝑞1 + 𝑟1 , where 0 ≤ 𝑟1 < 𝑏. If 𝑟1 ≠ 0, the g.c.d. (𝑎, 𝑏) is equal to
𝑏; 𝑖𝑓 𝑟1 ≠ 0, then (𝑎, 𝑏) = (𝑏, 𝑟1 ), so the step is repeated with 𝑏 and 𝑟1 in place of 𝑎 and
𝑏. After further repetitions, the last non-zero remainder obtained is the required g.c.d.
For example, for 𝑎 = 927 and 𝑏 = 643, write
 927 = 1 × 643 + 274,
643 = 2 × 274 + 95,
274 = 2 × 95 + 84,
95 = 1 × 84 + 11,
84 = 7 × 11 + 7,
11 = 1 × 7 + 4,
7 = 1 × 4 + 3,
4 = 1 × 3 + 1,
3 = 3 × 1,
and then 927, 643 = 643, 274 = 274, 95 = 95, 84 = 84, 11 = 11, 7 =
7, 4 = 4, 3 = 3,1 = 1.
EUCLIDEAN DOMAINS: Let 𝑹 be an integral domain, and let 𝑹∗ denote the set 𝑹\{𝟎} of
non-zero elements of 𝑹. An integer-valued function 𝝓: 𝑹∗ → 𝒁 defined on 𝑹∗ is said to
be a Euclidean function if it satisfies the following properties:—
(iv) 𝝓(𝒓) ≥ 𝟎 for all 𝒓 ∇ 𝑹∗ ;
(v) if 𝒙, 𝒚 ∇ 𝑹∗ satisfy 𝒙|𝒚 then 𝝓(𝒙) ≤ 𝝓(𝒚);
(vi) given 𝒙, 𝒚 ∇ 𝑹∗ , there exist 𝒒, 𝒓 ∇ 𝑹 such that
𝒙 = 𝒒𝒚 + 𝒓, where either 𝒓 = 𝟎 or 𝝓(𝒓) < 𝜙(𝑦).
A Euclidean domain is an integral domain on which is defined a Euclidean function.
EUCLIDEAN FIELD: An ordered field in which every positive element is a square. For
example, the field 𝑹 of real numbers is a Euclidean field. The field 𝑸 of rational numbers
is not a Euclidean field.
EUCLIDEAN GEOMETRY: The area of mathematics relating to the study of geometry
based on the definitions and axioms set out in Euclid’s book, the Elements. The
geometry of space described by the system of axioms first stated systematically in
the Elements of Euclid. The space of Euclidean geometry is usually described as a set of
objects of three kinds, called "points" , "lines" and "planes" ; the relations between them
are incidence, order, congruence and continuity. The parallel axiom (fifth postulate)
occupies a special place in the axiomatics of Euclidean geometry. The first sufficiently
precise axiomatization of Euclidean geometry was given by D. Hilbert. There are
modifications of Hilbert's axiom system as well as other versions of the axiomatics of
Euclidean geometry.
EUCLIDEAN ISOMETRIES: An isometry of 𝑬𝑵 is a map 𝑻 ∶ 𝑬𝑵 → 𝑬𝑵 that preserves the
Euclidean metric: 𝒅(𝑻(𝒙), 𝑻(𝒚) = 𝒅(𝒙, 𝒚) for all 𝒙, 𝒚 ∇ 𝑬𝑵 .
EUCLIDEAN PRIME NUMBER THEOREM: The set of prime numbers is infinite.
The Chebyshev theorems on prime numbers and the asymptotic law of the distribution
of prime numbers provide more precise information on the set of prime numbers in the
series of natural numbers.
EUCLIDEAN RING: An integral domain with an identity such that to each non-zero
element a of it corresponds a non-negative integer 𝒏(𝒂) satisfying the following
requirement: For any two elements 𝒂 and 𝒃 with 𝒃 ≠ 𝟎 one can find
elements 𝒒 and 𝒓 such that
𝑎 = 𝑏𝑞 + 𝑟,

where either 𝑟 = 0 or 𝑛(𝑟) < 𝑛(𝑏).

Every Euclidean ring is a principal ideal ring and hence a factorial ring; however, there
exist principal ideal rings that are not Euclidean. Euclidean rings include the ring of
integers (the absolute value |𝑎| plays the part of 𝑛(𝑎)), and also the ring of polynomials
in one variable over a field (𝑛(𝑎) is the degree of the polynomial). In any Euclidean ring
the Euclidean algorithm can be used to find the greatest common divisor of two
elements.
EUCLIDEAN SPACE (CARTESIAN SPACE): The number line 𝑹, plane 𝑹𝟐 and 3-
dimensional space 𝑹𝟑 can be generalized to n-dimensional space 𝑹𝒏 with co-ordinates
(𝑥1 , 𝑥2 , … , 𝑥𝑛 ) on which the operations of addition and multiplication by a scalar have
been extended in the obvious way. While 𝑹𝒏 is hard to visualize for 𝑛 > 3 it provides a
very powerful framework for multivariable analysis.

EUCLID–EULER THEOREM: Euclid–Euler theorem is a theorem in mathematics that

relates perfect numbers to Mersenne primes. An even positive integer is a perfect
number, that is, equals the sum of its proper divisors, if and only if it has the form
2p−1Mp where Mp is a Mersenne prime (i.e. a prime number of the form Mp = 2p − 1).

Suppose that 𝑚 ∇ 𝑁. If 2𝑚 − 1 is a prime, then the number 2𝑚 −1 (2𝑚 − 1) is an even

perfect number. Furthermore, there are no other even perfect numbers.
EUCLID NUMBERS: The perfect numbers which are even, for example 6 and 28 are
known as Euclid numbers.
EUCLID’S AXIOMS: The axioms Euclid set out in his famous text, the Elements, are:
 A straight line may be drawn from any point to any other point,
 A straight line segment can be extended indefinitely at either end,
 A circle may be described with any centre and any radius,
 Things that coincide with one another are equal to one another.
 The whole is greater than the part.
EUCLIDEAN SPACE OVER A FIELD: Let be a (commutative) field of characteristic not
two. A Euclidean space is a vector space 𝑿 over 𝑭 equipped with a symmetric bilinear
form 𝝈: 𝑿 × 𝑿 → 𝑭 satisfying 𝝈(𝒙, 𝒙) ≠ 𝟎 for all 𝒙 ∇ 𝑿, 𝒙 ≠ 𝟎. The elements of 𝑿 are
called points, and a set of points 𝒑 + 𝑭. 𝒗 (𝒑, 𝒗 ∇∇ 𝑿, 𝒗 ≠ 𝟎) is called a line. Let
𝑸 𝒙 = 𝝈(𝒙, 𝒙). Two points 𝒂, 𝒃 , (𝒄, 𝒅) are said to be congruent if and only if
𝑸 𝒂 − 𝒃 = 𝑸(𝒄 − 𝒅).

EULER ANGLES: The angles 𝜙, 𝜓 and 𝜃 that determine the position of one Cartesian
rectangular coordinate system 0𝑥𝑦𝑧 relative to another one 0𝑥′𝑦′𝑧′ with the same origin
and orientation. The Euler angles are regarded as the angles through which the former
must be successively rotated about the axes of the latter so that in the end the two
systems coincide.
EULER EQUATION: A linear ordinary differential equation of order 𝑛 of the form
𝑛
𝑑𝑖 𝑦
𝑎𝑖 𝑥 𝑖 = 𝑓(𝑥)
𝑑𝑥 𝑖
𝑖=0

where 𝑎𝑖 are constants and 𝑎𝑛 ≠ 0.

EULER FORMULAS: Formulas connecting the exponential and trigonometric functions:

𝑒 𝑖𝑧 = 𝑐𝑜𝑠𝑧 + 𝑖𝑠𝑖𝑛𝑧,

𝑒 𝑖𝑧 + 𝑒 −𝑖𝑧 𝑒 𝑖𝑧 − 𝑒 −𝑖𝑧
𝑐𝑜𝑠𝑧 = , 𝑠𝑖𝑛𝑧 = .
2 2𝑖

These hold for all values of the complex variable 𝑧. In particular, for a real value 𝑧 =
𝑥 the Euler formulas become
𝑒 𝑖𝑥 + 𝑒 −𝑖𝑥 𝑒 𝑖𝑥 − 𝑒 −𝑖𝑥
𝑐𝑜𝑠𝑥 = , 𝑠𝑖𝑛𝑥 = .
2 2𝑖

These formulas were published by L. Euler.

EULER-FROBENIUS POLYNOMIALS: The Euler–Frobenius polynomials 𝒑𝒎 (𝒙)

of degree 𝒎 − 𝟏 ≥ 𝟎 are characterized by the Frobenius reciprocal identity
𝟏
𝒙𝒎−𝟏 𝒑𝒎 = 𝒑𝒎 (𝒙)
𝒙
Thus, 𝒑𝒎 (𝒙) is invariant under the reflection of the indeterminate 𝑥.

EULER INTEGRALS: The integral

1

𝐵(𝑝, 𝑞) = 𝑥 𝑝−1 (1 − 𝑥)𝑞−1 𝑑𝑥 , 𝑝, 𝑞 > 0,

0

called the Euler integral of the first kind, or the beta-function, and

𝑥 𝑠−1 𝑒 −𝑥 𝑑𝑥
0

called the Euler integral of the second kind. (The latter converges for 𝑠 > 0 and is a
representation of the gamma-function.)

EULER, LEONHARD (1707–83): Leonord Euler was the most prolific of famous
mathematicians. He was born in Switzerland. He worked in a highly productive period
when the newly developed calculus was being extended in all directions at once, and he
made contributions to most areas of mathematics, pure and applied. Euler, more than
any other individual, was responsible for notation that is standard today. Among his
contributions to the language are the basic symbols, 𝜋, 𝑒 𝑎𝑛𝑑 𝑖, the summation notation
𝛴 and the standard function notation 𝑓(𝑥). His Introduction in analysis in infinitorum
was the most important mathematics text of the late eighteenth century.
EULER NUMBERS: The coefficients in the expansion

The recurrence formula for the Euler numbers ( in symbolic notation) has the
form
Thus, , the are positive and the are negative integers for
all ; , , , , and . The
Euler numbers are connected with the Bernoulli numbers by the formulas

EULERIAN CIRCUIT: An Eulerian circuit in a graph is a circuit that traverses every edge
of the graph. It follows from these definitions that any closed Eulerian trail is an
Eulerian circuit.
EULERIAN GRAPH: One area of graph theory is concerned with the possibility of
travelling around a graph, going along edges in such a way as to use every edge exactly
once. A connected graph is called Eulerian graph if there is a sequence 𝑣0 , 𝑒1 , 𝑣1 , … , 𝑒𝑘 , 𝑣𝑘
of alternately vertices and edges (where 𝑒𝑖 is an edge joining 𝑣𝑖−1 and 𝑣𝑖 ), with 𝑣0 = 𝑣𝑘
and with every edge of the graph occurring exactly once.
EULERIAN TRAIL: An Eulerian trail in a graph is a trail that traverses every edge of the
graph. Note that an Eulerian trail in a graph must traverse every edge of the graph
exactly once, since a trail traverses an edge of the graph at most once.

EULAR’S METHOD OF SUMMATION: In the series 𝑢𝑛 , if

𝑘+1 𝑘+1 𝑘+1

𝑠0 + 𝑠1 + ⋯ + 𝑠 2− 𝑘+1
1 2 𝑘+1 𝑘

𝑛
𝑠𝑛 = 𝑘=0 𝑢𝑘 converges to 𝑠 as 𝑘 → ∞, then 𝑢𝑛 is said to be summable by Eular’s
method, and we write 𝑢𝑛 = 𝑠(𝐸). The transformation thus obtained is called Eular’s
method of summation.

EULER PRODUCT: The infinite product

1
1− − 1,
𝑝𝑠
𝑝

where 𝑠 is a real number and 𝑝 runs through all prime numbers. This product converges
absolutely for all 𝑠 > 1. The analogous product for complex numbers 𝑠 = 𝜍 + 𝑖𝑡
converges absolutely for 𝜍 > 1 and defines in this domain the Riemann zeta-function
 1
𝜁(𝑠) = 1− − 1,
𝑝𝑠
𝑝

∞
1
= .
𝑛𝑠
𝑛=1

EULER STRAIGHT LINE: The straight line passing through the point 𝐻 of intersection of
the altitudes of a triangle, the point 𝑆 of intersection of its medians, and the centre 𝑂 of
the circle circumscribed to it. If the Euler line passes through a vertex of the triangle,
then the triangle is either isosceles or right-angled, or both right-angled and isosceles.
The segments of the Euler line satisfy the relation
𝑂𝐻: 𝑆𝐻 = 1: 2

1 1 1
EULER’S CONSTANT: Let 𝑎𝑛 = 1 + 2 + 3 + − − − + 𝑛 − ln 𝑛. This sequence has a limit

whose value is known as Euler’s constant, 𝛾; that is, 𝑎𝑛 → 𝛾. The value equals 0.577
215 66 to 8 decimal places. It is not known whether 𝛾 is rational or irrational.
𝑥
EULER’S EQUATIONS: The necessary condition for ∫𝑥 2 𝑓 𝑥, 𝑦, 𝑦′ 𝑑𝑥 to be an extremum
1

is that
𝜕𝑓 𝑑 𝜕𝑓
− =0
𝜕𝑦 𝑑𝑥 𝜕𝑦′
This is called Euler’s equation
EULER’S EQUATION OF MOTION ALONG A STREAMLINE: Consider an elementary
section of a stream tube. Let 𝛿𝑠 be the length of the stream tube element. Mass of the
fluid particle moving along a streamline in the positive direction is 𝜌 𝛿𝐴 𝛿𝑠. The force
acting on the element are of two types: (i) Body force and (ii) Surface force exerted due
to hydrostatics pressure on the end areas of the particle.

The body force is (𝜌𝐹𝑠 ) 𝛿𝐴𝛿𝑠. On the upstream face the pressure force is 𝜌𝛿𝐴 in the (+𝑠)
𝜕𝑝
direction face it is 𝜌 + 𝜕𝑠 𝛿𝐴 , 𝛿𝐴 acting in the (−𝑠) direction. The total force along the

path 𝛿𝑠 with tangential unit vector is given by

𝜕𝑝
𝜌𝐹𝑠 𝛿𝑠𝛿𝐴 + 𝜌𝛿𝐴 − 𝜌 + 𝛿𝐴
𝜕𝑠

𝜕𝐴
= 𝜌𝐹𝑠 𝛿𝑠𝛿𝐴 − 𝜕𝑠 𝛿𝑠𝛿𝐴.
 dq
The acceleration of the fluid flowing along 𝛿𝑠 is . By using Newton’s second law of
dt

motion the equation of momentum along the path is given by

dq 𝜕𝜌
𝜌𝛿𝑠𝛿𝐴 = 𝜌𝐹𝑠 𝛿𝑠𝛿𝐴 − 𝜕𝑠 𝛿𝑠𝛿𝐴
dt

dq 1 𝜕𝑝
Or = 𝐹𝑠 − 𝜌 𝜕𝑠 ,
dt

𝜕𝑝 𝜕𝑝 1 𝜕𝑝
Or + 𝑞 𝜕𝑠 = 𝐹𝑠 − 𝜌 𝜕𝑠 ,
𝜕𝑠

known as the Euler’s equation of motion for one- dimensional flow.

EULER’S FORMULA: The name given to the equation 𝑐𝑜𝑠 𝜃 + 𝑖 𝑠𝑖𝑛 𝜃 = 𝑒 𝑖𝜃 , a special
case of which gives 𝑒 𝑖𝜋 + 1 = 0.
EULER FORMULA ( DIFFERENTIAL GEOMETRY): A formula expressing the normal
curvature of a surface in a given direction l in terms of the principal
curvatures 𝑘1 and 𝑘2 :
𝑘𝑙 = 𝑘1 𝑐𝑜𝑠 2 𝜙 + 𝑘2 𝑠𝑖𝑛2 𝜙,

where ϕ is the angle between the direction 𝑙 and the principal direction corresponding
to the principal curvature 𝑘1 . This formula was established by L. Euler (1760).

EULER’S FUNCTION: For a positive integer 𝑛, let 𝜑(𝑛) be the number of positive integers
less than 𝑛 that are relatively prime to 𝑛. For example, 𝜑(12) = 4, since four numbers,
1, 5, 7 and 11, are relatively prime to 12. This function 𝜑, defined on the set of positive
integers, is Euler’s function. The arithmetic function 𝜙 whose value at 𝑛 is equal to the
number of positive integers not exceeding 𝑛 and relatively prime to 𝑛 . The Euler
function is a multiplicative arithmetic function, that is 𝜙(1) = 1 and 𝜙(𝑚𝑛) =
𝜙(𝑚)𝜙(𝑛) for (𝑚, 𝑛) = 1. The function 𝜙(𝑛) satisfies the relations

𝜙(𝑑) = 𝑛,
𝑑|𝑛

𝜙(𝑛) = 3𝜋 2 𝑥 2 + 𝑂(𝑥𝑙𝑛𝑥).
𝑛≤𝑥

EULER’S IDENTITY: The relation

 ∞ −1
1 1
= 1−
𝑛𝑠 𝑝𝑠
𝑛 =1 𝑝

where 𝑠 > 1 is an arbitrary real number and the product extends over all prime
numbers 𝑝. The Euler identity also holds for all complex numbers 𝑠 = 𝜍 + 𝑖𝑡 with 𝜍 > 1.

The Euler identity can be generalized in the form

∞ −1
1
𝑓(𝑛) = 1−
𝑓(𝑝)
𝑛=1 𝑝

which holds for every totally-multiplicative arithmetic function 𝑓(𝑛) > 0 for which the
∞
series 𝑛=1 𝑓(𝑛) is absolutely convergent.

EULER’S INTEGRALS- BATA AND GAMMA FUNCTIONS: We define the first and second
Eulerian integrals as

1
𝐵 𝐼, 𝑚 = 𝑥 𝑏−1 1 − 𝑥 𝑚 −1
𝑑 𝑥
0

and
∞
𝑇 𝑛 = 𝑒 −𝑥 𝑥 𝑛−1 𝑑 𝑥
0

and refer to them as the Beta and Gamma Functions respectively. These integrals
occupy a very important place in definite integrals and have got wide application in
Mechanics, Physics, Statistics and many other applied sciences.

We shall assume that all the quantities 𝐼, 𝑚, 𝑛 are positive but not necessarily integrals.

It will be noted that Beta Function is symmetrical in 𝐼 and 𝑚,

i.e.,
𝐵 𝐼, 𝑚 = 𝐵 𝑚, 𝑙

EULER’S METHOD: The simplest numerical method for solving differential equations. If
and an initial condition is known, 𝑦 = 𝑦0 when 𝑥 = 𝑥0 , then Euler’s method generates
a succession of approximations 𝑦𝑛+1 = 𝑦𝑛 + 𝑕𝑓 (𝑥𝑛 , 𝑦𝑛 ) where 𝑥𝑛 = 𝑥0 + 𝑛𝑕, 𝑛 =
1, 2, 3, …. This takes the known starting point, and moves along a straight line segment
with horizontal distance h in the direction of the tangent at (𝑥0 , 𝑦0 ). The process is
repeated from the new point (𝑥1 , 𝑦1 ) etc. If the step length 𝑕 is small enough, the
tangents are good approximations to the curve. The method provides a reasonably
accurate estimate.
EULER SERIES: The expression
1
,
𝑝
𝑝

where the sum extends over all prime number 𝑝. L. Euler (1748) showed that this series
diverges, thus providing another proof of the fact that the set of prime numbers is
infinite. The partial sums of the Euler series satisfy the asymptotic relation
1 1
= ln⁡
(ln 𝑥) + 𝐶 + 𝑂
𝑝 ln 𝑥
𝑝≤𝑥

where C=0.261497….
EULER’S THEOREM: Let 𝐺 be a connected planar graph drawn in the plane. If there are 𝑣
vertices, 𝑒 edges and 𝑓 faces, then 𝑣 – 𝑒 + 𝑓 = 2.
EVEN FUNCTION: The real function 𝑓 is an even function if 𝑓(– 𝑥) = 𝑓(𝑥) for all 𝑥 in the
domain of the function 𝑓. Thus the graph 𝑦 = 𝑓(𝑥) of an even function has the 𝑦-axis as
a line of symmetry. For example, 𝑓 is an even function when 𝑓(𝑥) is defined as any of the
2
following: 3, 𝑥 4 , 𝑥 6 + 4, 𝑥 4 + 1, 1/(𝑥 5 ), 𝑐𝑜𝑠 4 𝑥.
EVEN PERMUTATION: A rearrangement of the original ordering which can be obtained
by an even number of exchanges of pairs of elements.
EVENT: A subset of the sample space relating to an experiment. It is a set of outcomes of
an experiment; a subset of a sample space. For example, suppose that the sample space
for an experiment in which a coin is tossed three times is {HHH, HHT, HTH, HTT, THH,
THT, TTH, TTT}, and let A = {HHH, HHT, HTH, THH}. Then A is the event in which at
least two ‘heads’ are obtained. If, when the experiment is performed, the outcome is one
that belongs to A, then A is said to have occurred. The intersection A ∩ B of two events is
the event that can be described by saying that ‘both A and B occur’. The union A ∪ B of
two events is the event that ‘either A or B occurs’. Taking the sample space as the
universal set, the complement A′ of A is the event that ‘A does not occur’. The
probability Pr(A) of an event A is often of interest. The following laws hold:
  𝑃(𝐴 ∪ 𝐵) = 𝑃(𝐴) + 𝑃(𝐵) – 𝑃(𝐴 ∩ 𝐵).
 When 𝐴 and 𝐵 are mutually exclusive events, 𝑃(𝐴 ∪ 𝐵) = 𝑃(𝐴) + 𝑃(𝐵).
 When 𝐴 and 𝐵 are independent events, 𝑃(𝐴 ∩ 𝐵) = 𝑃(𝐴) 𝑃(𝐵).
 𝑃(𝐴′) = 1 – 𝑃(𝐴).
EVOLVENT: An evolvent is a region in the plane that is isometric to a given region on a
developable surface. Example: The evolvent of the side surface of a cone cut along a
generator is a planar sector. The approximate construction of an evolvent can be
achieved graphically by means of descriptive geometry.
EXACT DIFFERENTIAL: If 𝑧 = 𝑓(𝑥, 𝑦) is a function of two independent variables then
𝜕𝑧 𝜕𝑧
𝑑𝑧 = 𝑑𝑥 + 𝑑𝑦
𝜕𝑥 𝜕𝑥
is the exact differential. For a function of more than two variables the exact differential
will have similar partial derivative terms for each of its independent variables.
EXACT DIFFERENTIAL EQUATION: An equation in which the exact differential of a
function is equal to zero. If 𝑧 = 𝑓(𝑥, 𝑦) then
𝜕𝑧 𝜕𝑧
𝑑𝑥 + 𝑑𝑦 = 0
𝜕𝑥 𝜕𝑥
is an exact differential equation.
EXACT DIVISOR: An exact divisor of a number is a factor of a given number. e.g. 7 is an
exact divisor of 35.
EXACT HOMOMORPHISM: Let 𝑹 be a unital ring, let 𝑭, 𝑮 and 𝑯 be 𝑹-modules, and let
𝒑: 𝑭 → 𝑮 and 𝒒: 𝑮 → 𝑯 be 𝑹-module homomorphisms. The sequence 𝑭(𝒑) −→
𝑮(𝒒) −→ 𝑯 of modules and homomorphisms is said to be exact at 𝑮 if and only if
𝒊𝒎𝒂𝒈𝒆 (𝒑: 𝑭 → 𝑮) = 𝒌𝒆𝒓(𝒒: 𝑮 → 𝑯). A sequence of modules and homomorphisms is
said to be exact if it is exact at each module occurring in the sequence (so that the image
of each homomorphism is the kernel of the succeeding homomorphism). A
monomorphism is an injective homomorphism. An epimorphism is a surjective
homomorphism. An isomorphism is a bijective homomorphism.
EXAMPLE: A particular instance of a generalized statement. A counter-example will
disprove a generalized claim, but examples do not provide proof. EXCELLENT RING: A
commutative Noetherian ring satisfying the three axioms stated below. It is known that
a geometric ring possesses several qualitative properties not inherent in arbitrary
Noetherian rings. The concept of an excellent ring makes it possible to take the most
important properties of geometric rings axiomatically into account.

Axioms of an excellent ring 𝐴.

A1: The ring 𝐴 is a universal chain ring. (A ring 𝐴 is said to be a chain ring if for any two
prime ideals 𝑃 ≠ 𝑃′ of it the lengths of any two saturated chains 𝑃 = 𝑃0 ⊆ 𝑃1 ⊆
𝑃2 … … … ⊆ 𝑃𝑛 = 𝑃′ of prime ideals are the same. A ring 𝐴 is said to be a universal
chain ring if any polynomial ring 𝐴[𝑇1 , 𝑇2 , … … … , 𝑇𝑘 ] is a chain ring.)

A2: The formal fibres of 𝐴 are geometrically regular, i.e. for any prime ideal 𝑃 ⊂ 𝐴 and
any homomorphism from 𝐴 into a field 𝐾, the ring 𝐴𝑝 ⨂𝑘 𝐾 is regular. Here 𝐴𝑝 is the
completion of the local ring 𝐴𝑝 .

A3: For any integral finite 𝐴 -algebra there is a non-zero element 𝑏 ∇ 𝐵 such that the
ring of fractions, 𝐵[𝑏 −1 ], is regular.

Excellent rings possess the following properties:

1) For an excellent ring 𝐴, the set of regular (normal) points of the scheme 𝑆𝑝𝑒𝑐 𝐴 is
open.

2) If an excellent local ring 𝐴 is reduced (normal or equi-dimensional), then so is the

completion 𝐴.

3) The integral closure of an excellent ring 𝐴 in a finite extension of the field of fractions
of 𝐴 is a finite 𝐴 -algebra.

4) If a ring 𝐴 is excellent, then any 𝐴 -algebra of finite type is also an excellent ring.

Two important examples of excellent rings are the complete local rings (or analytic
rings) and the Dedekind rings with field of fractions of characteristic zero. Therefore,
the class of excellent rings is sufficiently large and contains, in particular, all algebras of
finite type over a field or over the ring 𝑍 of integers.

The excellence of a ring 𝐴 is closely connected with the possibility of resolution of

singularities of the scheme 𝑆𝑝𝑒𝑐 𝐴.

EXCEPTIONAL ANALYTIC SET: An analytic set in a complex space 𝑿 for which there
exists an analytic mapping 𝒇: 𝑿 → 𝒀 such that 𝒇(𝑨) → 𝒀 is a point in the complex
space 𝒀, while 𝒇: 𝑿 − 𝑨 → 𝒀 − 𝒚 is an analytic isomorphism. The modification 𝒇 is
called a contraction of 𝑨 to 𝒚.
EXCLUSIVE DISJUNCTION: The proposition 𝐴 ∨ ˙𝐵, obtained from two
propositions 𝐴 and 𝐵 using the exclusive disjunction ∨ ˙, is taken to be true if 𝐴 is true
and 𝐵 is false, or if 𝐴 is false and 𝐵 is true. In the remaining cases it is taken to be false.
Thus, the exclusive disjunction can be expressed in terms of the ordinary (non-
exclusive) disjunction by the formula
𝐴 ∨ ˙𝐵 ⇔ (𝐴 ∨ 𝐵)&¬(𝐴 ∧ 𝐵).

EXISTENCE THEOREM: If 𝑘 𝑠 and 𝜏(𝑠) are continuous functions of a real variable

𝑠(𝑠 ≥ 0) then there exists a space curve for which 𝑘 is the curvature, 𝜏 is the torsion,
and 𝑠 is the arc length measured from some suitable base point.
EXHAUSTIVE EVENTS: A set of events in statistics whose union is the whole probability
space, or a set of sets whose union is the universal set under consideration.
exp: It is the abbreviation for exponential function. (exp x is also written as ex).
EXPANSION OF SOME BASIC FUNCTIONS:
𝑥2 𝑥3
 𝑒𝑥 = 1 + 𝑥 + + +−−−−
2! 3!
𝑥2 𝑥3
 𝑒 −𝑥 = 1 − 𝑥 + − +−−−−
2! 3!
𝑥2 𝑥3
 𝑎 𝑥 = 1 + 𝑥 𝑙𝑜𝑔 𝑎 + (𝑙𝑜𝑔 𝑎)2 + (𝑙𝑜𝑔 𝑎)3 + − − − −
2! 3!
𝑥2 𝑥3
 𝑎−𝑥 = 1 − 𝑥 𝑙𝑜𝑔 𝑎 + 2!
(𝑙𝑜𝑔 𝑎)2 − 3!
(𝑙𝑜𝑔 𝑎)3 + − − − −
𝑥2 𝑥3
 𝑙𝑜𝑔 1 + 𝑥 = 𝑥 − + −−−−−
2! 3!
𝑥2 𝑥3
 𝑙𝑜𝑔 1 − 𝑥 = − 𝑥 + + +−−−−
2! 3!

EXPECTED VALUE: The expected value 𝐸(𝑋) of a random variable 𝑋 is a value that gives
the mean value of the distribution, and is defined as follows. For a discrete random
variable 𝑋, 𝐸(𝑋) = 𝛴 𝑝𝑖 𝑥𝑖 , where 𝑝𝑖 = 𝑃 (𝑋 = 𝑥𝑖 ). For a continuous random variable
𝑋,
∞

𝐸 𝑥 = 𝑥𝑓 𝑥 𝑑𝑥
−∞

where 𝑓 is the probability density function of 𝑋. The following laws hold:

 𝐸(𝑎𝑋 + 𝑏𝑌) = 𝑎𝐸(𝑋) + 𝑏𝐸(𝑌).
  When 𝑋 and 𝑌 are independent, 𝐸(𝑋𝑌) = 𝐸(𝑋) 𝐸(𝑌).
EXPERIMENT: An action that has measurable or quantifiable results.
EXPLICIT FORMULA: A formula for the nth term of a sequence, written as an expression
of n.
EXPLICIT FUNCTION: If the dependent variable y is expressed in the form 𝑦 = 𝑓(𝑥)
then 𝑦 is an explicit function of 𝑥. So 𝑦 = 6𝑥 + 15 is explicit but 6𝑥 – 𝑦 + 15 = 0 is
not, though it can be rearranged to be explicit.
EXPONENTIAL DISTRIBUTION: The continuous probability distribution with probability
density function 𝑓 given by 𝑓(𝑥) = 𝜆 𝑒𝑥𝑝(– 𝜆 𝑥), where 𝜆 is a positive parameter, and
𝑥 ≥ 0. It has mean 1/𝜆 and variance 1/𝜆2 . The time between events that occur
randomly but at a constant rate has an exponential distribution. The distribution is
skewed to the right.
EXPONENTIAL EQUATION: An equation in which the variable occurs in an exponent.
EXPONENTIAL FUNCTION: The function 𝑓 such that 𝑓(𝑥) = 𝑒 𝑥 , or 𝑒𝑥𝑝 𝑥, for all 𝑥 in 𝑅.
The two notations arise from different approaches described below, but are used
interchangeably. A function of the form y = Abx, where A is a non-zero constant and b >
0 and b ≠ 1, is also called exponential function.
EXPONENTIAL GROWTH: Growth that can be modelled by an exponential function

EXPONENTIAL LAW: A term with a variety of meanings

 A process which changes by a fixed factor over a fixed time period: exponential
growth or decay.

 A property of the exponential function, such as 𝑒𝑥𝑝(𝑥 + 𝑦) = 𝑒𝑥𝑝(𝑥)𝑒𝑥𝑝(𝑦) or of

a power, such as 𝑥 𝑦 𝑧
= 𝑥 𝑦𝑧 .
 A property of sets of maps between sets, such
as 𝑀𝑎𝑝(𝑋 × 𝑌, 𝑍) ↔ 𝑀𝑎𝑝(𝑋, 𝑀𝑎𝑝(𝑌, 𝑍)).
 A property of sets of continuous maps between topological spaces.
Ext: It is the abbreviation for Ext functor.

ext: It is the abbreviation for exterior.

EXTENDABLE FUNCTION: A function 𝑓 ∶ 𝑋 → 𝑌 is extendable, 𝑓 ∇ 𝐸𝑥𝑡(𝑋, 𝑌 ), provided
there exists a connectivity function 𝐹 ∶ 𝑋 × [0, 1] → 𝑌 such that 𝑓 (𝑥) = 𝐹 (𝑥, 0) for
every 𝑥 ∇ 𝑋. It is easy to see that 𝐶(𝑋, 𝑌 ) ⊂ 𝐸𝑥𝑡(𝑋, 𝑌 ) ⊂ 𝐶𝑜𝑛𝑛(𝑋, 𝑌 ) ⊂ 𝐷(𝑋, 𝑌 ) for
arbitrary topological spaces, where 𝐶(𝑋, 𝑌 ) stands for the class of all continuous
functions from 𝑋 into 𝑌.
EXTENDED COMPLEX PLANE: It is the set of complex numbers with a point at infinity.
The set can be denoted by 𝑪∞ and can be thought of as a Riemann sphere by means of a
stereographic projection. If a sphere is placed so that a point 𝑆 on the sphere is touching
the complex plane at the origin, then 𝑆 corresponds to the point (0,0) on the complex
plane, which is the complex number 𝑧 = 0. All other points on the sphere, except 𝑁
which is diametrically opposite to 𝑆 on the sphere, are in a one-to-one correspondence
with points on the complex plane through the stereographic projection, and therefore
with a unique complex number. The point 𝑁 is identified with the point at infinity, with
corresponding complex number ∞.
EXTENDED REAL NUMBERS: It is the set of real numbers, with the positive and negative
infinite cardinals.
EXTENSION OF A FIELD: A field extension 𝑲 is a field containing a given field k as a
subfield. The notation 𝑲/𝒌 means that 𝑲 is an extension of the field 𝒌. In this case, 𝑲 is
sometimes called an overfield of the field 𝒌.
Let 𝐾/𝑘 and 𝐿/𝑘 be two extensions of a field 𝑘. An isomorphism of fields 𝜑: 𝐾 → 𝐿 is
called an isomorphism of extensions if 𝜑 is the identity on 𝑘. If an isomorphism of
extensions exists, then the extensions are said to be isomorphic. If 𝐾 = 𝐿, φ is called an
automorphism of the extension 𝐾/𝑘. The set of all automorphisms of an extension
forms a group, 𝐴𝑢𝑡(𝐾/𝑘). If 𝐾/𝑘 is a Galois extension, this group is denoted
by 𝐺𝑎𝑙(𝐾/𝑘) and is called the Galois group of the field 𝐾 over 𝑘, or the Galois group of
the extension 𝐾/𝑘. An extension is called Abelian if its Galois group is Abelian.
EXTENSION OF A MODULE: Any module 𝑿 containing the given module 𝑨 as a
submodule. Usually one fixes a quotient module 𝑿/𝑨, that is, an extension of the
module by the module is an exact sequence
𝟎→𝑨→𝑿→𝑩→𝟎

Such a module 𝑋 always exists (for example, the direct sum of 𝐴 and 𝐵), but need not be
uniquely determined by 𝐴 and 𝐵.

EXTENSION OF AN OPERATOR: A linear operator whose graph contains the graph of the
given linear operator. When the operator 𝑩 is an extension of a given operator 𝑨, one
writes 𝑨 ⊂ 𝑩. The usual problems in the theory of extensions are as follows: extend an
operator maximally while preserving a specific property, or study extensions of an
operator having various additional properties.

For example, let 𝐴 be a given isometric operator on a Hilbert space 𝐻 with domain of
definition 𝐷(𝐴) ⊂ 𝐻 and range of values 𝑅(𝐴) ⊂ 𝐻; then the isometric extensions
of 𝐴 are in one-to-one correspondence with the isometric mappings from 𝐻+ =
𝐷(𝐴)⊥ to 𝐻− = 𝑅(𝐴)⊥ . In particular, 𝐴 has unitary extensions if the dimensions
of 𝐻+and 𝐻− coincide.

EXTENSION OF CAUCHY’S INTEGRAL FORMULA TO MULTIPLY CONNECTED REGIONS:

If 𝑓(𝑧) is a analytic in a ring shaped region bounded by two closed curves 𝐶1 𝑎𝑛𝑑 𝐶2 and
a is a point in the region between 𝐶1 𝑎𝑛𝑑 𝐶2 .

1 𝑓(𝑧) 1 𝑓(𝑧)
𝑓 𝑎 = ∫𝐶 𝑑𝑧 − ∫𝐶 𝑑𝑧
2𝜋𝑖 2 𝑧 − 𝑎 2𝜋𝑖 1 𝑧 − 𝑎

Where 𝐶2 is outer curve.

EXTENSION OF CAUCHY’S THEOREM COMPLEX ANALYSIS): Suppose 𝑓(𝑧) is analytic in

a simply connected domain 𝐷. Then the integral along any rectifiable curve in 𝐷 joining
any given points of 𝐷 is the same i.e., it does not depend upon the curve joining the two
points.

EXTENSION OF TOPOLOGICAL SPACE: A topological space 𝒀 in which the given

topological space 𝑿 is an everywhere-dense set. If 𝒀 is a compact space, then it is called
a compact extension, and if 𝒀 is a Hausdorff space, it is called a Hausdorff extension.
EXTERIOR NORMAL TO A CONVEX SURFACE: A vector perpendicular to a supporting
hyperplane and directed towards that half-space defined by the supporting hyperplane
which does not contain the points of the surface.
EXTRAPOLATION: An extension of the function beyond the boundary of its domain of
definition, in which the extended (as a rule) function belongs to a given class. The
extrapolation of functions is usually performed by means of formulas using information
about the behaviour of the function at a finite collection of points (interpolation nodes)
in its domain of definition. The concept of interpolation of functions is used as the
opposite concept to that of extrapolation and consists in constructively (possibly,
approximately) re-establishing the values of functions in their domains of definition.
EXTREMAL FIELD: A domain in the (𝑛 + 1) −dimensional space of the variables
𝑥, 𝑦1 , 𝑦2 , … … … , 𝑦𝑛 covered without intersections by an 𝑛-parameter family of extremals
of the functional
⌌𝐵⌍

𝐽= 𝐹 𝑥, 𝑦1 , 𝑦2 , … … … , 𝑦𝑛 , 𝑦 ′ , … … … , 𝑦𝑛′ 𝑑𝑥
⌌𝐴⌍

where 𝐴 and 𝐵 are the initial and final points through which the extremals of the family
pass. One must distinguish between proper (or general) and central extremal fields. A
proper extremal field corresponds to the case when the extremals of the family are
transversal to some surface

𝜑 𝑥, 𝑦1 , 𝑦2 , … … … , 𝑦𝑛 = 0.

For a proper extremal field the point 𝐴 (or 𝐵) in (1) belongs to the surface (2) and the
condition (3) is satisfied at it.

A central extremal field corresponds to the case when the extremals of the family
emanate from one point lying outside the field, for example, from a common initial
point 𝐴. The slope of an extremal field is the vector-function
𝑢 𝑥, 𝑦 = (𝑢1 𝑥, 𝑦 , … … … , 𝑢𝑛 𝑥, 𝑦 associating with every point
𝑥, 𝑦 = (𝑥, 𝑦1 , 𝑦2 , … … … , 𝑦𝑛 ) of the field the vector 𝑦 𝑥 = (𝑦1′ 𝑥 , … … … . 𝑦𝑛′ 𝑥 ).

EXTREMALLY-DISCONNECTED SPACE: A space in which the closure of every open set is

open. In a regular extremally-disconnected space there are no convergent sequences
without repeated terms. Therefore, among the metric spaces only the discrete ones are
extremally disconnected. Nevertheless, extremally-disconnected spaces are fairly
widespread: Every Tikonov space can be represented as the image under a perfect
irreducible mapping of some extremally-disconnected Tikonov space. This means that
extremal disconnectedness is not preserved by perfect mappings. However, the image
of an extremally-disconnected space under a continuous open mapping is an
extremally-disconnected space.
EXTREMAL METRIC, METHOD OF THE: One of the fundamental methods in geometric
function theory, closely connected with differential geometry and topology. The method
of the extremal metric is based on the relations between the length of curves belonging
to specific homotopy classes and the areas of the domains filled out by them. Here these
curves and areas are computed in a special metric corresponding to the peculiarities of
the extremal problem under investigation.

The method of the extremal metric has various forms. The original one is Grötzsch's
strip method, which is an essential refinement of the arguments connecting length and
area, operating with the characteristic conformal invariants of doubly-connected
domains and quadrangles. Using his strip method, H. Grötzsch obtained a number of
classical results in the theory of conformal and quasi-conformal mappings

EXTREMAL PROBLEM: The problem of finding the extrema of functions or functionals

by choosing parameters or functions (controls) from various classes defined by various
conditions (phase, differential, integral, etc.) imposed on the classes over which the
specific value, function or functional is to attain a minimum or maximum.
EXTREMAL PROPERTIES OF FUNCTIONS: Properties of individual functions that
distinguish them as solutions of some extremal problems. The majority of the special
functions arising in mathematical analysis can be characterized by some extremal
property. This applies, for example, to extremal properties of polynomials: the
classical Laguerre polynomials, the Legendre polynomials, the Chebychev polynomials,
the Hermite polynomials. The classical polynomials are usually solutions of various
extremal problems arising not infrequently in remote domains of analysis. For example,
the Chebychev polynomials are extremal in the problem of an inequality for the
derivatives of polynomials. The same can also be said of other special functions. Many of
them are eigen functions of differential operators, that is, they are solutions of
some isoperimetric problem. In this context the best known special functions are
connected in some way with the existence of an invariant structure, when they are eigen
functions of the Laplace–Beltrami equation, which is shift-invariant. This holds for the
trigonometric polynomials, the spherical functions, the cylinder functions, etc. The
majority of extremal properties of functions can be stated in the form of some exact
inequality.
EXTREME POINT OF A CONVEX SET: A point x in a convex set C is called an extreme
point if x cannot be expressed as a convex combination of any two distinct points 𝑥1 and
𝑥2 in C.

Mathematically, a point is a extreme point of a convex set if there do not exist other
points 𝑥1 ,𝑥2 (𝑥1 ≠ 𝑥2 ) in the set , such that
 𝑥 = 𝜆𝑐1 + 1 − 𝜆 𝑥2 , 0 ≤ 𝜆 ≤ 1

We say that 𝑥 is an extreme point of a convex set 𝑆 if whenever 𝑥 = 𝜃𝑦 + (1 − 𝜃)𝑧,

for 𝑦, 𝑧 ∇ 𝑆, 0 < 𝜃 < 1, then 𝑥 = 𝑦 = 𝑧.

EXTREME VALUE OF A FUNCTION: It is the maximum or minimum value of a function. A

function 𝑓: 𝐴 → 𝐵 is numeric if 𝐵 is a set of numbers. For a numeric function, it's
possible to compare its values at different points 𝑎 ∇ 𝐴. Extreme is a value which, in
some sense is either maximum or minimum. If 𝑓(𝑎) exceeds all other values of 𝑓 then
we say it's a global extreme, maximum. If it's only larger than values of 𝑓 for points
near 𝑎, the maximum is local.

EXTREME VALUE THEOREM : Extreme value theorem states that if a real-valued

function 𝑓 is continuous in the closed and bounded interval [𝑎, 𝑏], then 𝑓 must attain
a maximum and a minimum, each at least once. That is, there exist numbers 𝑐 and 𝑑 in
[𝑎, 𝑏] such that:

A related theorem is the boundedness theorem which states that a continuous

function 𝑓 in the closed interval [𝑎, 𝑏] is bounded on that interval. That is, there exist
real numbers 𝑚 and 𝑀 such that:

EXTREMUM: A point at which a function has a turning point, i.e. at least a local
maximum or minimum. An extremum is a point where a function attains a maximum or
minimum. This article will consider only functions that are continuous and
differentiable. A global maximum is the point where a function attains its highest value.
A local maximum is a point where the value of the function is higher than the
surrounding points. Similar definitions apply to minimum points. Both local maximum
and local minimum points can be found by determining where the curve has a
horizontal tangent, which means that the derivative is zero at that point. If the first
derivative is zero and the second derivative is positive, then the curve is concave up,
and the point is a minimum.
F
FACTORIAL NOTATION: The notation, 𝑛!, used to represent the product of the first n
natural numbers.

FACTORIAL RING: A ring with unique decomposition into factors. More precisely, a
factorial ring 𝑨 is an integral domain in which one can find a system of irreducible
elements 𝑷 such that every non-zero element 𝒂 ∇ 𝑨 admits a unique representation

𝒂=𝒖 𝒑𝒏(𝒑)
𝒑∇𝑷

where 𝑢 is invertible and the non-negative integral exponents 𝑛(𝑝) are non-zero for
only a finite number of elements 𝑝 ∇ 𝑃. Here an element is called irreducible
in 𝐴 if 𝑝 = 𝑢𝑣 implies that either 𝑢 or 𝑣 is invertible in 𝐴, and 𝑝 is not invertible in 𝐴.
FACTOR THEOREM: Suppose that 𝑓(𝑥) represents a polynomial in 𝑥. The factor theorem
says that, if 𝑓 𝑎 = 0, then (𝑥 − 𝑎) is one of the factors of 𝑓(𝑥).
FACTOR SPACE: If 𝑋 is a vector space and 𝑀 is a subspace of 𝑋, then 𝑋/𝑀 is a vector
space, called the factor space (or the quotient space) of 𝑋 with respect to 𝑀.

F-ALGEBRA: A real vector space that is simultaneously a lattice is called a vector

lattice (or Riesz space) whenever 𝒙 ≤ 𝒚 ( ≤ is the lattice order) implies 𝒙 + 𝒛 ≤ 𝒚 +
𝒛 for all 𝒛 ∇ 𝑨 and 𝜶𝒙 ≤ 𝜶𝒚 for all positive real numbers 𝜶. If 𝑨 is also
an algebra and 𝒛𝒙 ≤ 𝒛𝒚 and 𝒙𝒛 ≤ 𝒚𝒛 for all 𝒛 ∇ 𝑨+, the positive cone of 𝑨, then 𝑨 is
called an 𝒍-algebra (a lattice-ordered algebra, Riesz algebra). A Riesz algebra 𝑨 is called
an 𝒇-algebra whenever
𝐢𝐧𝐟 𝒙, 𝒚 = 𝟎 ⇒ 𝐢𝐧𝐟 𝒛𝒙, 𝒚 = 𝐢𝐧𝐟 𝒙𝒛, 𝒚 = 𝟎∀𝒛 ∇ 𝑨+

This notion was introduced by G. Birkhoff and R.S. Pierce in 1956.

An important example of an 𝑓-algebra is 𝐴 = 𝐶(𝑋), the space of continuous functions on

some topological space 𝑋. Other examples are spaces of Baire functions, measurable
functions and essentially bounded functions.

FAMILY OF CURVES: A set of similar curves which are of the same form and
distinguished by the values taken by one or more parameters in their general equation.
In particular, where the solution of a differential equation is obtained, the general
solution will involve one or more constants of integration, giving rise to a family of
curves. A particular member of the family may be identified as the required solution if
boundary conditions are known.

FAMILY OF SURFACE: An equation of the form

𝐹 𝑥, 𝑦, 𝑧, 𝑎 = 0
Where ‘𝑎’ is a constant represents a family of surface. If ‘𝑎’ can take all real values i.e. if
‘𝑎’ is parameter, then equation represents one parameter family of surface having ‘𝑎’ as
parameter. By assigning different values to ‘𝑎’ we shall get different surface of this
family of surface.
Similarly the equation of the form 𝐹 𝑥, 𝑦, 𝑧, 𝑎, 𝑏 =0, where ‘𝑎’ and ‘𝑏’ are parameters
represents the family of surface having ‘𝑎’ and ‘𝑏’ as two parameters.
FARKAS’ LEMMA: Let 𝐶 be a closed convex cone in 𝑅 𝑚 and let 𝑏 be a vector in 𝑅 𝑚 .
Suppose that 𝑏 ∈ 𝐶. Then there exists a linear functional 𝜙: 𝑅 𝑚 → 𝑅 such that
𝜙(𝑣) ≥ 0 for all 𝑣 ∇ 𝐶 and 𝜙(𝑏) < 0.
FATOU LEMMA: If a sequence (𝑓𝑛 ) of µ −summable non-negative functions is such that:

 ∫ 𝑓𝑛 𝑑µ ≤ 𝐶 for all 𝑛;
𝑎. 𝑒
 𝑓𝑛 𝑓.,
→

then 𝑓 is µ-summable and ∫ 𝑓 𝑑µ ≤ 𝐶. 𝐿1 (𝑋) is a Banach space.

F-DISTRIBUTION: The F-distribution is a continuous random variable distribution that
is frequently used in statistical inference. It is used to test the hypothesis that two
normally distributed random variables have the same variance, and in regression to test
the relationship between an explanatory variable and the dependent variable. The
distribution is skewed to the right. Tables relating to the distribution are available.
FEASIBLE SOLUTION OF AN LPP: The feasible solution if a LPP is that set of values
which satisfies both the constraints and the non – negative restrictions of a LPP. A
feasible solution is a set of values for the choice variables in a linear programming
problem that satisfies the constraints of the problem. Some properties of CPF Solutions
are:

Property 1 If there is exactly one optimal solution, it must be a corner-point

feasible (CPF) solution.
Property 2 If there are multiple optimal solutions ≥ 2 and a bounded feasible
region, optimal solutions must be adjacent corner-point feasible (CPF)
solution.
Property 3 There are a finite number of corner-point feasible (CPF) solutions.
Property 4 If a corner-point feasible (CPF) solution. has no adjacent corner-point
feasible (CPF) solution. that are better, then such a corner-point
feasible (CPF) solution. is guaranteed to be an optimal solution.

 Properties also hold for Basic feasible solution.

 Each Basic feasible solution has m basic variables, and the rest are non-basic.
 The number of non-basic variables equals n + m, where m is number of surplus
variables.
 The basic solution is the augmented corner-point solution whose n defining
equations are indicated by the non-basic variables.
 A Basic feasible solution is a basic solution where all m basic variables are non-
negative.
 A Basic feasible solution is said to be degenerate if any of the m variables equals
zero.

FEJÉR’S THEOREM: Let 𝑓 ∇ 𝐶𝑃[−𝜋, 𝜋]. Then

 1 𝜋 1 𝑛 𝑘
𝐹𝑛 𝑥 = 2𝜋 ∫−𝜋 𝑓 𝑡 𝐾𝑛 𝑥 − 𝑡 𝑑𝑡, 𝑤𝑕𝑒𝑟𝑒 𝐾𝑛 𝑡 = 𝑛 +1 𝑘=0 𝑚 =−𝑘 𝑒 𝑖𝑚𝑡 is the Fejér

kernel.

FEKETE–SZEGŐ INEQUALITY : Fekete–Szegő inequality is an inequality for the

coefficients of univalent analytic functions found byFekete and Szegő (1933), related to
the Bieberbach conjecture. Finding similar estimates for other classes of functions is
called the Fekete–Szegő problem.

The Fekete–Szegő inequality states that if

is a univalent analytic function on the unit disk and 0 ≤ λ < 1, then

FERMAT, PIERRE DE FERMAT: Pierre de Fermat (1601 to 1665) was French

mathematician who developed number theory, worked on ideas that later became
known as calculus, and corresponded with Pascal on probability theory. He is
remembered chiefly for his work in the theory of numbers, including Fermat’s Little
Theorem and what is known as Fermat’s Last Theorem. His work on tangents was an
acknowledged inspiration to Newton in the latter’s development of the calculus. Fermat
introduced coordinates as a means of studying curves. Professionally he was a judge in
Toulouse, and to mathematicians he is the ‘Prince of Amateurs’.
𝑟
FERMAT PRIME: A prime of the form 22 + 1 is called as a Fermat prime as were
informed by Fermat. At present, the only known primes of this form are those given by r
= 0, 1, 2, 3 and 4.
FERMAT’S LAST THEOREM: Fermat’s last theorem states that there is no solution to the
equation 𝑎𝑛 + 𝑏 𝑛 = 𝑐 𝑛 where 𝑎, 𝑏, 𝑐, and 𝑛 are all positive integers, and 𝑛 > 2. Fermat
wrote in the margin of a book that he had a proof of this, but as he never repeated the
claim it is likely that he realized the incompleteness of his supposed proof. Much
research was done over centuries until a proof was completed in 1995 by Andrew
Wiles.
FERMAT’S LITTLE THEOREM: Let 𝑝 be a prime, and let 𝑎 be an integer not divisible by
𝑝. Then 𝑎𝑝 – 1 ≡ 1 (𝑚𝑜𝑑 𝑝).

FERMAT'S THEOREM ON SUMS OF TWO SQUARES: An odd prime 𝑝 is expressible as

with 𝑥 and 𝑦 integers, if and only if

For example, the primes 5, 13, 17, 29 and 37 are all congruent to 1 modulo 4,

and they can be expressed as sums of two squares in the following ways:

FERRARI METHOD: The Ferrari method is a method for reducing the solution of an
equation of degree 4 over the complex numbers to the solution of one cubic and two
quadratic equations; it was discovered by L. Ferrari.

The Ferrari method for the equation

𝑦 4 + 𝑎𝑦 3 + 𝑏𝑦 2 + 𝑐𝑦 + 𝑑 = 0

consists in the following. By the substitution 𝑦 = 𝑥 − 𝑎/4 the given equation can be
reduced to

𝑥 4 + 𝑝𝑥 2 + 𝑞𝑥 + 𝑟 = 0, (1)

which contains no term in 𝑥 3 . If one introduces an auxiliary parameter 𝛼, the left-hand

side of (1) can be written as

𝑝2
𝑥 4 + 𝑝𝑥 2 + 𝑞𝑥 + 𝑟 = (𝑥 2 + 𝑝2 + 𝛼)2 − [2𝛼𝑥 2 − 𝑞𝑥 + (𝛼 2 + 𝑝𝛼 + − 𝑟)]. (2)
4

One then chooses a value of α such that the quadratic trinomial in square brackets is a
perfect square. For this the discriminant of the quadratic trinomial must vanish. This
gives a cubic equation for 𝛼,

𝑝2
𝑞 2 − 4 ⋅ 2𝛼(𝛼 2 + 𝑝𝛼 + − 𝑟) = 0.
4

Let α0 be one of the roots of this equation. For 𝛼 = 𝛼0 the polynomial in square brackets
in (2) has one double root,

𝑞
𝑥0 = ,
4𝛼0
which leads to the equation

(𝑥 2 + 𝑝2 + 𝛼0 )2 − 2𝛼0 (𝑥 − 𝑥0 )2 = 0.

This equation of degree 4 splits into two quadratic equations. The roots of these
equations are also the roots of (1).

FIBONACCI (1170–1250): An Italian merchant by the name of Leonardo of Pisa, he was

one of those who introduced the Hindu–Arabic number system to
Europe. He strongly advocated this system in Liber abaci, published in 1202, which also
contained problems including one that gives rise to the Fibonacci numbers.
FIBONACCI SEQUENCE: The first two numbers of the Fibonacci sequence are 1; every
other number is the sum of the two numbers that immediately precede it. Therefore, the
first 14 numbers in the sequence are: 1, 1, 2, 3, 5, 8, 13, 21, 34, 55, 89, 144, 233, 377,
610, 987, 1597,….. This sequence has many interesting properties. For instance, the
sequence consisting of the ratios of one Fibonacci number to the previous one, , has the
limit τ, the golden ratio.
FIELD: A field is a set of elements having two binary operations defined on it with
following properties:
—It is an Abelian group with respect to one operation called addition (with an identity
element designated 0).
—It is also an Abelian group with respect to another operation called multiplication.
—The distributive property holds: 𝑎 𝑏 + 𝑐 = 𝑎𝑏 + 𝑎𝑐.
Roughly speaking, S is a field if addition, subtraction, multiplication and division (except
by zero) are all possible in S. We normally use the letter K for a general field. For
example, the real numbers are an example of a field, with addition and multiplication
defined in the traditional manner. The concept can also be generalized to other types of
objects.
There are many other fields, including some finite fields. For example, for each prime
number p, there is a field 𝔽𝑝 = 0,1,2, ⋯ ⋯ , 𝑝 − 1 with 𝑝 elements, where addition and
multiplication are carried out modulo 𝑝. Thus, in 𝔽7 , we have 5 + 4 = 2 ,5 × 4 =
6 𝑎𝑛𝑑 5−1 = 3 𝑏𝑒𝑐𝑎𝑢𝑠𝑒 5 × 3 = 1.The smallest such field 𝔽2 has just two elements 0 and
1 where 1 + 1 = 0.

Various other familiar properties of numbers, such as

 0𝛼 = 0, −𝛼 𝛽 = − 𝛼𝛽 , −𝛼 −𝛽 = 𝛼𝛽, −1 𝛼 = −𝛼 , 𝑓𝑜𝑟 𝑎𝑙𝑙 𝛼, 𝛽 𝜖 𝑆

can be proved from the axioms.

However, occasionally you need to be very careful. For example in 𝔽2 we have 1+1=0,
and so it is not possible to divide by 2 in this field.

FIELD EXTENSIONS : Let 𝐾 be a field. An extension 𝐿: 𝐾 of 𝐾 is an embedding of 𝐾 in

some larger field 𝐿.
FIELD OF EXTREMALS: A central field (Proper) is called a field of extremals, if it is
formed by a family of extremals.
FILTERS: A filter on a set 𝑀 is a non-empty family 𝐹 ∇ 𝑃(𝑀) such that:
 𝐴, 𝐵 ∇ 𝐹 ⇒ 𝐴 ∩ 𝐵 ∇ 𝐹;
 𝐴 ∇ 𝐹 𝑎𝑛𝑑 𝐴 ⊂ 𝐵 ⇒ 𝐵 ∇ 𝐹;
 𝜑 ∈ 𝐹.
A filter is called free if ∩ 𝐹 = 𝜑.
FINITE ABELIAN GROUPS: An Abelian group 𝐺 of order 𝑝, where 𝑝 is a prime number, is
a direct product of cyclic subgroups 𝑍1 , . . . , 𝑍𝑟 , : 𝐺 = 𝑍1 × … × 𝑍𝑟 . If 𝑍𝑖 is of order 𝑝𝑛 𝑖 ,
then 𝑛 = 𝑛1 + − − − + 𝑛𝑟 , and we can assume that 𝑛𝑖 ≥ 𝑛𝑖+1 . A direct product
decomposition of 𝐺, as above, is not unique, but 𝑛1 , − − −, 𝑛𝑟 are determined uniquely
by 𝐺. The system {𝑝𝑛 1 , − − −, 𝑝𝑛 𝑟 } or {𝑛1 , − − −, 𝑛𝑟 } is called the system of invariants
(or type) of 𝐺, and a system of generators {𝑧1 , , − − −, 𝑧𝑟 } of 𝑍1 × … × 𝑍𝑟 is called a basis
of 𝐺. An Abelian group of type (𝑝, 𝑝, − − −, 𝑝) is called an elementary Abelian group.
FINITE AND 𝝈 −FINITE MEASURES: A measure µ is finite if µ(𝐴) < ∞ for all 𝐴 ⊂ 𝑋. A
measure µ is 𝜍 −finite if 𝑋 is a union of countable number of sets 𝑋𝑘 , such that the
restriction of µ to each 𝑋𝑘 is finite.
For a measure µ on a 𝜍 −algebra R, we have the following results

1. If 𝐴, 𝐵 ∇ 𝑅 with 𝐴 ⊆ 𝐵, then µ(𝐴) ≤ µ(𝐵)

2. If 𝐴, 𝐵 ∇ 𝑅 with 𝐴 ⊆ 𝐵 and µ(𝐵) < ∞, then µ(𝐵 ∖ 𝐴) = µ(𝐵) − µ(𝐴);

3. If (𝐴𝑛 ) is a sequence in R, with 𝐴1 ⊆ 𝐴2 ⊆ 𝐴3 ⊆ ⋯. Then lim𝑛→∞ µ 𝐴𝑛 =

µ(⋃𝑛 𝐴𝑛 )

4. If (𝐴𝑛 ) is a sequence in R, with 𝐴1 ⊇ 𝐴2 ⊇ 𝐴3 ⊇ ⋯.If µ 𝐴𝑛 < ∞ for some 𝑚,

then
 lim𝑛→∞ µ 𝐴𝑛 = µ(⋂𝑛 𝐴𝑛 ).

Any measure µ′ on a semiring 𝑆 is uniquely extended to a measure µ on the generated

ring 𝑅(𝑆). If the initial measure was 𝜍 −additive, then the extension is 𝜍 −additive as
well.

FINITE DEFORMATION: The deformation in which displacements 𝓊 together with their

derivatives are no longer small is called finite deformation.

FINITE-DIMENSIONAL VECTOR SPACE: A vector space is said to be finite-dimensional if

every element of it can be expressed as a linear combination of a finite set of linearly
independent vectors.
FINITE ELEMENT METHOD: A numerical method of solving partial differential equations
with boundary conditions by considering a series of approximations which satisfy the
differential equation and boundary conditions within a small region
FINITE FOURIER COSINE TRANSFORM: Let f(t) be a function defined on 0, ℓ and
satisfying Dirichlet’s conditions on 0, ℓ . The finite Fourier cosine transform of
𝑓 𝑡 , 𝑜 < 𝑡 < 𝑙 𝑖𝑠 𝑑𝑒𝑓𝑖𝑛𝑒𝑑 𝑎𝑠

𝓵
𝑐𝑜𝑠𝜋𝑆𝑡
𝐹𝑐 𝑓 𝑡 = 𝑓 𝑡 𝑑𝑡; 𝑠 𝜖 𝑁
𝟎 ℓ

FINITE FOURIER SINE TRANSFORM: Let 𝑓(𝑡) be a function defined on 0, ℓ and

satisfying Dirichlet’s conditions on 0, ℓ . The finite Fourier sine transform of
𝑓 𝑡 , 0 < 𝑡 < 𝑙 is defined as

𝓵
𝜋𝑆𝑡
𝐹𝑠 𝑓 𝑡 = 𝑓 𝑡 𝑠𝑖𝑛 𝑑𝑡; 𝑠 𝜖 𝑁
𝟎 ℓ

FINITE INTERSECTION PROPERTY: A space has the finite intersection property when
every family of closed subsets such that all finite sub collections have non-empty
intersections means that the entire family also has a non-empty intersection. It can be
shown that this property is equivalent to the space being compact.
FINITELY GENERATED ABELIAN GROUPS: The theory of finitely generated Abelian
groups, i.e., Abelian groups generated by a finite number of elements, is as old as that of
finite Abelian groups. The direct product of infinite cyclic groups is called a free Abelian
group. A finitely generated Abelian group G is the direct product of a finite Abelian
group and a free Abelian group. The finite factor is the torsion subgroup of G. The free
factor of the group G is not necessarily unique; however, the number of infinite cyclic
factors of the free factor is uniquely determined and is called the rank of G. Two finitely
generated Abelian groups are isomorphic if they have isomorphic maximal torsion
subgroups and the same rank.
FINITELY GENERATED IDEAL: An ideal 𝑰 of the ring 𝑹 is said to be finitely generated if
there exists a finite subset of 𝑹 which generates the ideal 𝑰. Every ideal of the ring 𝒁 of
integers is generated by some non-negative integer 𝒏.
FINITE MATHEMATICS: The branch of mathematics concerned with the study of
properties of structures of finite character, that arise both within mathematics and in
applications. Among these structures one has, e.g., finite groups, finite graphs, and also
certain mathematical models of information processing, finite automata, Turing
machines, etc.. Sometimes the subject of finite mathematics is assumed to extend to
arbitrary discrete structures, and this leads to discrete mathematics, identifying the
latter with finite mathematics. Certain algebraic systems, infinite graphs, definite forms
of computing schemes, cellular automata, etc., can be regarded as belonging to this area.
The term discrete analysis sometimes serves as a synonym for the concepts of "finite
mathematics" and "discrete mathematics" .
FINITE RANK OPERATOR: Let 𝑋 and 𝑌 be normed spaces, 𝑇 ∇ 𝐵 𝑋, 𝑌 is a finite rank
operator if 𝐼𝑚 𝑇 is a finite dimensional subspace of 𝑌.
FINITE SEQUENCE: A sequence that has a last term.

FINSLER MANIFOLD: A Finsler manifold allows the definition of distance but does not
require the concept of angle; it is an analytic manifold in which each tangent space is
equipped with a norm, || · ||, in a manner which varies smoothly from point to point.
This norm can be extended to a metric, defining the length of a curve; but it cannot in
general be used to define an inner product. Any Riemannian manifold is a Finsler
manifold.

FIRST-COUNTABLE SPACES: The space (𝑋, 𝑇) is first countable if each point has a
countable base of neighborhoods, i.e., for each point 𝑝 in 𝑋 there is a countable
collection of open sets such that each neighborhood of 𝑝 contains at least one of them,
for all 𝑝 ∇ 𝑋 : ∃ {𝑂𝑛 | 𝑛 ∇ 𝐍, 𝑂𝑛 ∇ 𝑇 for all 𝑛}, such that for all 𝑈, 𝑝 ∇ 𝑈 ∇
𝑇 : ∃ 𝑂𝑛 ⊂ 𝑈 . All metric spaces are first countable (the countable bases of
neighborhoods are the balls of radius 1/𝑛).

FIRST FUNDAMENTAL FORM OR METRIC: First fundamental form for an embedding or

immersion is the pullback of the metric tensor. Let 𝑟 = 𝑟(𝑢, 𝑣) be the equation of a
surface and let 𝐸 = 𝑟12 = 𝑟1 . 𝑟1 , 𝐹 = 𝑟1 . 𝑟2 and 𝐺 = 𝑟22 = 𝑟2 . 𝑟2 . The quadratic
differential form

𝐸𝑑𝑢² + 2 𝐹𝑑𝑢𝑑𝑣 + 𝐺𝑑𝑣²

is 𝑑𝑢, 𝑑𝑣 is called metric or first fundamental form of the surface and the quantities
𝐸, 𝐹, 𝐺 are called first fundamental coefficients or first order fundamental magnitudes.
Since 𝐸, 𝐹 and 𝐺 are functions of 𝑢, 𝑣 the quantities will generally vary form point to
point on the surface.

FIRST ISOMORPHISM THEOREM: Let 𝐺 be a group, let 𝐻 be a subgroup of 𝐺, and let 𝑁 be

a normal subgroup of 𝐺. Then 𝐻𝑁/ 𝑁 ≅ 𝐻 /𝑁 ∩ 𝐻 .
FIRST-ORDER DIFFERENTIAL EQUATION: A differential equation containing only the
differential coefficients of the first order. For example,
𝑑𝑦 𝑑𝑦
= 7𝑦, 𝑦 + 3𝑥 = 0
𝑑𝑥 𝑑𝑥
etc are first-order differential equations.
FISHER, RONALD AYLMER (1890–1962): Ronald Aylmer Fisher was a British statistician
who established methods of designing experiments and analysing results that have been
extensively used ever since. He developed the t-test and the use of contingency tables,
and is responsible for the method known as ANOVA.
FIXED FIELD OF A GROUP: Let 𝐿 be a field, and let 𝐺 be a group of automorphisms of 𝐿.
The fixed field of 𝐺 is the subfield 𝐾 of 𝐿 defined by 𝐾 = {𝑎 ∇ 𝐿 ∶ 𝜍(𝑎) = 𝑎 for
all 𝜍 ∇ 𝐺}.

FIXED POINTS: The points which coincide with their transformation are called invariant
or fixed points of the transformation that is to say, fixed points of a transformation
𝑤 = 𝑓(𝑧) are obtained by the equation 𝑧 = 𝑓 𝑧 .
FIXED POINT THEOREM: Any theorem which gives conditions, under which a mapping
must have a fixed point, is known as Fixed point theorem. They have been important in
the development of mathematical economics.
FLOOR: For a real number 𝑟, its floor value [𝑟] is defined as the largest integer not
greater than 𝑟. Thus [5] = [5.1] = 5 and [−5] = −5 while [−5.1] = −6. A more recent
notation ⌊𝑟⌋ allows to distinguish with another frequently used function - ceiling -
⌈𝑟⌉ which is defined as the smallest integer not less than 𝑟.
𝑑𝑥
FLOQUET'S THEOREM: Let 𝑑𝑡 = 𝐴 𝑡 𝑥 be a linear first order differential equation,

where 𝑥(𝑡) is a column vector of length 𝑛 and 𝐴 𝑡 an 𝑛 × 𝑛 periodic matrix with

period 𝑇. Let 𝜑(𝑡) be a fundamental matrix solution of this differential equation. Then,
for all 𝑡 ∇ ℝ,

FLOQUET THEORY: A theory concerning the structure of the space of solutions, and the
properties of solutions, of a linear system of differential equations with periodic
coefficients
𝒙′ = 𝑨 𝒕 𝒙; 𝒕 ∇ 𝑹, 𝒙 ∇ 𝑹𝒏 (𝟏)

the matrix 𝐴(𝑡) is periodic in 𝑡 with period 𝜔 > 0 and is summable on every compact
interval in 𝑅.

Every fundamental matrix 𝑋 of the system (1) has a representation

𝑋 𝑡 = 𝐹 𝑡 exp⁡
(𝑡𝐾)

called the Floquet representation, where 𝐹 𝑡 is some 𝜔-periodic matrix and 𝐾 is some
constant matrix. There is a basis 𝑥1 , 𝑥2 , … … … , 𝑥𝑛 of the space of solutions of (1) such
that 𝐾 has Jordan form in this basis; this basis can be represented in the form

𝑥𝑖 = 𝜓1𝑖 exp( 𝛼𝑖 𝑡 , … … … , 𝜓𝑛𝑖 exp( 𝛼𝑖 𝑡))

where 𝜓𝑘𝑖 are polynomials in 𝑡 with 𝜔-periodic coefficients, and the 𝛼𝑖 are the
characteristic exponents of the system (1). Every component of a solution of (1) is a
linear combination of functions of the form (of the Floquet solutions) 𝜓𝑘𝑖 exp( 𝛼𝑖 𝑡). In
the case when all the characteristic exponents are distinct, the 𝜓𝑘𝑖 are simply 𝜔-
periodic functions. The matrices 𝐹(𝑡) and 𝐾 in the representation (2) are, generally
speaking, complex valued. If one restricts oneself just to the real case, then 𝐹(𝑡) does
not have to be 𝜔-periodic, but must be 2𝜔-periodic.

FLUID: By definition, a fluid cannot withstand any tendency for applied forces to deform
it in such a way that volume is left unchanged. Such deformation may be resisted, but
not prevented, by a fluid. Ordinary gases and liquids are fluids.
FOCI “Foci” is the plural of “focus.”
FOCUS (1) A parabola is the set of points that are the same distance from a fixed point
(the focus) and a fixed line (the directrix). The focus, or focal point, is important
because starlight striking a parabolically shaped telescope mirror will be reflected back
to the focus.
(2) An ellipse is the set points such that the sum of the distances to two fixed points is a
constant. The two points are called foci. Planetary orbits are shaped like ellipses, with
the sun at one focus.
FOCUS OF A CURVE :A point 𝑭 lying in the plane of the second-order curve such that the
ratio of the distance of any point of the curve from 𝑭 to its distance from a given line
(the directrix) is equal to a constant (the eccentricity). The foci of a second-order curve
can be defined as the points of intersection of the tangents to that curve from
the circular points of the plane. This definition can also be extended to algebraic curves
of order 𝒏.
FOL: It is the abbreviation for first-order logic.
FOREST: A forest is an acyclic graph.
FORT’S THEOREM: If f is discontinuous at the points of a dense set, show that the set of
points x, where f′(x) exists, is of the first category.
FORWARD DIFFERENCE: If {(𝑥𝑖 , 𝑓𝑖 )}, 𝑖 = 0, 1, 2, … is a given set of function values with
𝑥𝑖+1 = 𝑥𝑖 + 𝑕, 𝑓𝑖 = 𝑓(𝑥𝑖 ) then the forward difference at 𝑓𝑖 is defined by 𝑓𝑖+1 – 𝑓𝑖 =
𝑓 𝑥𝑖+1 − 𝑓(𝑥𝑖 ).
FOURIER INTEGRAL FORMULA: Let 𝑓 𝑥 be a function satisfying Dirichlet’s Conditions
∞
and is absolutely integrable in −∞ ∞ i.e. ∫−∞ f x dx converges,

∞
Then 𝑓 𝑥 = ∫0 𝐹 𝑠 cos 𝑠 𝑥 + 𝐺 𝑠 sin 𝑠 𝑥 𝑑𝑠

∞ ∞
Where 𝜋 𝐹 𝑠 = ∫−∞ f w cos s w dw, π G s = ∫−∞ f w sin s w dw
Also then 𝑓 𝑥 can be re-written as

∞ ∞ ∞
1 1
𝑓 𝑥 = 𝑓 𝑤 cos 𝑠 𝑤 𝑑𝑤 cos 𝑠 𝑥 + 𝑓 𝑤 sin 𝑠 𝑤 𝑑𝑤 sin 𝑠 𝑥 𝑑𝑠
0 𝜋 −∞ 𝜋 −∞
𝑠=∞ 𝑤=∞
1
= 𝑓 𝑤 cos 𝑠 𝑤 cos 𝑠 𝑥 + sin 𝑠 𝑤 sin 𝑠𝑥 𝑑𝑤 𝑑𝑠
𝜋 𝑠=0 𝑤=−∞

𝑠=∞ 𝑤=∞
1
= 𝑓 𝑤 cos 𝑠 𝑥 − 𝑤 𝑑𝑤 𝑑𝑠
𝜋 𝑠=0 𝑤=−∞

𝑠=∞ 𝑤=∞
1
= 𝑓 𝑤 cos 𝑠 𝑥 − 𝑤 𝑑𝑤 𝑑𝑠
2𝜋 𝑠= −∞ 𝑤=−∞

The representation of 𝑓 𝑥 in Fourier integral formula hold if 𝑓 𝑥 is continuous at 𝑥 .

𝑓 𝑥+0 + 𝑓 𝑥−𝑜
However if 𝑥 is a point of discontinuity, then 𝑓 𝑥 is replaced by as in case
2

of Fourier series. As for Fourier series, the above conditions are necessary but not
sufficient.

∞
Also ∫0 𝐹 𝑠 cos 𝑠 𝑥 + 𝐺 𝑠 sin 𝑠 𝑥 𝑑𝑠 is known as Fourier integral expansion of 𝑓 𝑥
or simply Fourier integral.

Different forms of Fourier integral formula:

1 ∞ ∞
(A) 𝑓 𝑥 = ∫−∞ 𝑓 𝑤 ∫0 cos 𝑠 𝑥 − 𝑤 𝑑𝑠 𝑑𝑤 is known as General Form
𝜋
2 ∞ ∞
(B) 𝑓 𝑥 = ∫0 sin 𝑠 𝑥 ∫0 𝑓 𝑤 sin 𝑠 𝑤 𝑑𝑤 𝑑𝑠 is known as Fourier sine
𝜋

integral formula if 𝑓 𝑥 is an odd function.

2 ∞ ∞
(C) 𝑓 𝑥 = 𝜋 ∫0 cos 𝑠 𝑥 ∫0 𝑓 𝑤 cos 𝑠 𝑤 𝑑𝑤 𝑑𝑠 is known as Fourier cosine

Integral formula if 𝑓 𝑥 is an even function

1 ∞ ∞
(D)𝑓 𝑥 = 2𝜋 ∫−∞ e𝒊 𝒔 𝒙 ∫−∞ f w e−i s w dw ds or
1 ∞ ∞
𝑓 𝑥 = 2𝜋 ∫−∞ e−𝒊 𝒔 𝒙 ∫−∞ f w ei s w dw ds is known as Complex or

exponential Form of Fourier integral formula.

FOURIER, JEAN-BAPTISTE JOSEPH (1768 TO 1830): Jean-Baptiste Joseph Fourier was a

French mathematician who studied differential equations of heat conduction, and
developed the concept now known as Fourier series. These so-called Fourier series are
of immense importance in physics, engineering and other disciplines, as well as being of
great mathematical interest.
FOURIER SERIES: Any periodic function can be expressed as a series involving sines and
cosines, known as a Fourier series. Assume that units are chosen so that the period of
the function is 2𝜋. Then:
𝑎0
𝑓 𝑥 = + 𝑎1 𝑐𝑜𝑠 𝑥 + 𝑏1 sin 𝑥 + 𝑎2 𝑐𝑜𝑠 2𝑥 + 𝑏2 sin 2𝑥 + 𝑎3 𝑐𝑜𝑠 3𝑥 + 𝑏3 sin 3𝑥 + −
2
−−−
where the Fourier coefficients are found from these integrals:
𝜋 𝜋
1 1
𝑎𝑛 = 𝑓 𝑥 cos 𝑛𝑥 𝑑𝑥 , 𝑏𝑛 = 𝑓 𝑥 sin 𝑛𝑥 𝑑𝑥
𝜋 𝜋
−𝜋 −𝜋

The derivation of theses coefficients requires some knowledge about the orthogonality
properties of the trigonometric functions. The main properties that we need are given
by the following integrals,

𝐿 𝒏𝝅
∫−𝐿 𝑠𝑖𝑛 𝑳
𝒙𝑑𝑥 = 0

𝐿 𝒏𝝅
∫−𝐿 𝑐𝑜𝑠 𝑳
𝒙𝑑𝑥 = 0, 𝑚 ≠ 0

𝐿 𝑚𝝅 𝒏𝝅
∫−𝐿 𝑐𝑜𝑠 𝑳
𝒙𝑐𝑜𝑠 𝑳
𝒙𝑑𝑥 = 0, 𝑖𝑓 𝑚 ≠ 𝑛

𝐿 𝑚𝝅 𝒏𝝅
∫−𝐿 𝑠𝑖𝑛 𝑳
𝒙𝑠𝑖𝑛 𝑳
𝒙𝑑𝑥 = 0, 𝑖𝑓 𝑚 ≠ 𝑛

𝐿 𝑚𝝅 𝒏𝝅
∫−𝐿 𝑐𝑜𝑠 𝑳
𝒙𝑠𝑖𝑛 𝑳
𝒙𝑑𝑥 = 0,

𝐿 𝒏𝝅 𝐿 𝒏𝝅
∫−𝐿 𝑠𝑖𝑛2 𝑳
𝒙𝑑𝑥 = ∫−𝐿 𝑐𝑜𝑠 2 𝑳
𝒙𝑑𝑥 = 𝐿

FOURIER TRANSFORM: The integral transform

∞

𝐹 𝑦 = 𝑓(𝑥)𝑒 𝑖𝑦𝑥 𝑑𝑥
−∞

is called the Fourier transformation. The function 𝐹 is said to be the Fourier transform
of the function 𝑓. Let 𝑓 𝑡 be a function defined and piecewise continuous on −∞ ∞
∞
and is absolutely convergent on −∞ ∞ i.e. ∫−∞ 𝒇 t dt converges. Then Fourier
transform of 𝑓 𝑡 , denoted by
 ∞
𝐹 𝑓 𝑡 and is defined as 𝐹 𝑓 𝑡 = ∫−∞ 𝒇 𝒕 ei s t dt = G s say

and the function 𝒇 𝒕 is inverse Fourier transform of 𝐺 𝑠 written as

𝑓 𝑡 = 𝐹 −1 𝐺 𝑠 ,

and defined as

∞
1
𝑓 𝑡 = G s
2𝜋 −∞

This is also known as Inversion formula for Fourier transform.

Relation between Laplace and Fourier transform:

Let 𝑓 𝑡 be a function defined as

𝑒 −𝑎 𝑡 𝑔 𝑡 , 𝑡 > 0
𝑓 𝑡 =
0, 𝑡<0

The Fourier transform of 𝑓 𝑡 = 𝐹 𝑓 𝑡

∞
= 𝒇 𝒕 ei b t dt
−∞

𝟎 ∞
ibt
= 𝒇 𝒕 e dt + 𝒇 𝒕 ei b t dt
−∞ −∞

𝟎 ∞
= 𝟎. ei b t dt + e−at g t ei b t dt
−∞ 𝟎

∞
= 0+ e− −a−i b t g t dt
𝟎

∞
= 𝑒 −𝑠 𝑡 g t dt, where s = a − i b
𝟎

=𝐿 𝑔 𝑡 = 𝐿𝑎𝑝𝑙𝑎𝑐𝑒 𝑡𝑟𝑎𝑛𝑠𝑓𝑜𝑟𝑚 𝑜𝑓 𝑔 𝑡 .
FOUR SQUARES THEOREM: Any positive integer can be expressed as the sum of not
more than four positive integers.
FOUR-VERTEX THEOREM: Four-vertex theorem states that the curvature function of a
simple, closed, smooth plane curve has at least four local extrema (specifically, at least
two local maxima and at least two local minima). The name of the theorem derives from
the convention of calling an extreme point of the curvature function a vertex.
FRACTION: A fraction is a rational number expressed in the form 𝑛/𝑑, where 𝑛 is
designated the numerator and 𝑑 the denominator. The slash between them is known as
a solidus when the fraction is expressed as 𝑛/𝑑. If 𝑛/𝑑 < 1, then 𝑛/𝑑 is known as a
proper fraction. Otherwise, it is an improper fraction. If 𝑛 and 𝑑 are relatively prime,
then 𝑛/𝑑 is said to be in lowest terms. To get a fraction in lowest terms, simply divide
the numerator and the denominator by their greatest common divisor.
FRACTIONAL IDEAL: A fractional ideal of 𝑅 is a subset of the field 𝐾 of fractions of 𝑅 that
is of the form 𝑐𝐼, where 𝑐 is a non-zero element of 𝐾 and 𝐼 is an ideal of 𝑅. A fractional
ideal of 𝑅 is not an ideal of 𝑅 unless it is contained in 𝑅. If 𝑅 is not itself a field, then the
field 𝐾 of fractions of 𝐾 will contain fractional ideals of 𝑅 that are not ideals of 𝑅. Indeed,
given any element 𝑐 of 𝐾 \ 𝑅, the subset 𝑐𝑅 of 𝐾 is a fractional ideal of 𝐾, where
𝑐𝑅 = {𝑐𝑟 ∶ 𝑟 ∇ 𝑅}, but it is not an ideal of 𝑅. However any fractional ideal of 𝑅 that is
contained in 𝑅 is an ideal of 𝑅.
FREDHOLM EQUATION, NUMERICAL METHODS: Approximation methods for solving
Fredholm integral equations of the second kind, amounting to performing a finite
number of operations on numbers.

Let

𝜑 𝑥 − 𝜆 ∫ 𝐾 𝑥, 𝑠 𝜙 𝑠 𝑑𝑠 = 𝑓(𝑥) (1)

be a Fredholm integral equation of the second kind, where 𝜆 is a complex

number, 𝑓(𝑥) is a known vector function, 𝜙(𝑥) is an unknown vector function, 𝐾(𝑥, 𝑠) is
the kernel of equation (1), and 𝐷 is a domain in some 𝑚-dimensional Euclidean space.
Below it is assumed that 𝜆 does not belong to the spectrum of the integral operator with
kernel 𝐾 (that is, for the given 𝜆 equation (1) has a unique solution in some class of
functions corresponding to the degree of smoothness of 𝐾). The expression (1)
naturally includes the case of a system of Fredholm equations.
FREDHOLM OPERATOR: A linear normally-solvable operator 𝑩 acting on a Banach
space 𝑬 with index 𝝌𝑩 equal to zero. The classic example of a Fredholm operator is an
operator of the form
𝑩=𝑰+𝑻

where is the identity and is a completely-continuous operator on .

FREDHOLM'S THEOREM: If 𝑀 is a matrix, then the orthogonal complement of the row

space of 𝑀 is the null space of 𝑀:

Similarly, the orthogonal complement of the column space of 𝑀 is the null space of the
adjoint:

FREE AND BOUND VARIABLES: Any function whose values are propositions is called a
propositional function. ∀𝑥 and ∃𝑥 can be regarded as operators that transform any
propositional function 𝐹(𝑥) into the propositions ∀𝑥𝐹(𝑥) and ∃𝑥𝐹(𝑥), respectively.
∀𝑥 and ∃𝑥 are called quantifiers; the former is called the universal quantifier and the
latter the existential quantifier. F(x) is transformed into ∀𝑥𝐹(𝑥) or ∃𝑥𝐹(𝑥) just as a
1
function 𝑓(𝑥) is transformed into the definite integral ∫0 𝑓(𝑥)𝑑𝑥; the resultant
propositions ∀𝑥𝐹(𝑥) and ∃𝑥𝐹(𝑥) are no longer functions of 𝑥. The variable 𝑥 in∀𝑥𝐹(𝑥)
and in ∃𝑥𝐹(𝑥) is called a bound variable, and the variable 𝑥 in 𝐹(𝑥), when it is not
bound by ∀𝑥 or ∃𝑥, is called a free variable.
FREE BASIS OF A MODULE: Let 𝑴 be a module over a unital commutative ring 𝑹.
Elements 𝒙𝟏 , 𝒙𝟐 , . . . , 𝒙𝒌 of 𝑴 are said to constitute a free basis of 𝑴 if these elements are
distinct, and if the 𝑹-module 𝑴 is freely generated by the set {𝒙𝟏 , 𝒙𝟐 , . . . , 𝒙𝒌 }.
FREE BOOLEAN ALGEBRA: A Boolean algebra with a system of generators such that
every mapping from this system into a Boolean algebra can be extended to a
homomorphism. Every Boolean algebra is isomorphic to a quotient algebra of some free
Boolean algebra.

For any cardinal number 𝛼 there is a unique (up to an isomorphism) free Boolean
algebra with 𝛼 generators. A finite Boolean algebra is free if and only if its number of
𝑛
elements is of the form 2𝑛 , where 𝑛 is the number of generators. Such a free Boolean
algebra is realized as the algebra of Boolean functions of 𝑛 variables. A countable free
Boolean algebra is isomorphic to the algebra of open-closed subsets of the Cantor set.
Every set of pairwise-disjoint elements of a free Boolean algebra is finite or countable.

An infinite free Boolean algebra cannot be complete. On the other hand, the cardinality
of any infinite complete Boolean algebra is the least upper bound of the cardinalities of
its free subalgebras.

FREE GROUP: A group 𝐹 with a system 𝑋 of generating elements such that any mapping
from 𝑋 into an arbitrary group 𝐺 can be extended to a homomorphism from 𝐹 into 𝐺.
Such a system 𝑋 is called a system of free generators; its cardinality is called the rank
of . The set 𝑋 is also called an alphabet. The elements of 𝐹 are words over the
alphabet 𝑋.
In other terms, Groups have generators, such elements that all other elements of a
group could be obtained from generators and their inverses using the group operation.
A group is said to be free if no relation exists between its generators other than between
an element and its inverse. The additive group of integers is free with a single generator
1. The multiplicative group of all positive rational numbers has prime numbers as its
generators. From the Fundamental Theorem of Arithmetic representation of integers in
the form 𝑝𝑎 𝑞 𝑏 . . . 𝑟 𝑐 where 𝑝, 𝑞, . . . , 𝑟 are all different primes, is unique. Therefore the
group is free.
FREE MODULE: A module over a unital commutative ring is said to be free if there exists
some subset of the module which freely generates the module.

FREE MODULE OF FINITE RANK: A module 𝑀 over an integral domain 𝑅 is said to be a

free module of finite rank if there exist elements 𝑏1 , 𝑏2 , . . . , 𝑏𝑘 ∇ 𝑀 that constitute a free
basis for 𝑀. These elements constitute a free basis if and only if, given any element 𝑚 of
𝑀, there exist uniquely-determined elements 𝑟1 , 𝑟2 , . . . , 𝑟𝑘 of R such that 𝑚 = 𝑟1 𝑏1 +
𝑟2 𝑏2 + · · · + 𝑟𝑘 𝑏𝑘 .

FREELY GENERATED MODULE: Let 𝑀 be a module over a unital commutative ring 𝑅,

and let 𝑋 be a subset of 𝑀. The module 𝑀 is said to be freely generated by the set 𝑋 if
the following conditions are satisfied:

(i) the elements of 𝑋 are linearly independent over the ring 𝑅;

(ii) the module 𝑀 is generated by the subset 𝑋.
FRITZ JOHN STATIONARY-POINT NECESSARY OPTIMALITY THEOREM: Let x be a local
solution of the problem: Minimize 𝜃 x , subject to x ∇ X 0 , g(x) ≤ 0, where X 0 is open set
in Rn and 𝜃𝑖 : X 0 → R, g: X 0 → Rm are differentiable at x. Then there exists an r0 ∇ R and
an r0 ∇ Rm such that x, r0 , r solves the problem r0 ∆θ x) + r ∆g(x = 0 subject to
g(x) ≤ 0, r g x = 0,(r0 , rw ) ≥ 0 where

W = i g i x = 0, and g i is concave at x

FRITZ JOHN STATIONARY- POINT PROBLEM (FJP): The Fritz John stationary- point
problem means to find x ∇ X 0 , r0 ∇ Rm if they exist, such that ∆x ∅ x, r0 , r = 0,
∆r ∅ x, r0 , r ≤ 0, r∆r ∅ x, r0 , r = 0, r0 , r ≥ 0 and ∅ x, r0 , r = r0 ∅ x) + rg(x or
equivalently, r0 ∆θ x) + r ∆g(x = 0, g(x) ≤ 0, r g x = 0, (r0 , r) ≥ 0.

Frob: It is the abbreviation for Frobenius endomorphism.

FRONTIER (BOUNDARY) POINTS: The set of points which are members of the closure of
a set and also of the closure of its complement. For example, for an open interval (𝑎, 𝑏)
or closed interval [𝑎, 𝑏], the values 𝑎 and 𝑏 are the frontier or boundary points.
FRUSTUM: A frustum is a portion of a cone or a pyramid bounded by two parallel
planes. A frustum of a right-circular cone is the part between two parallel planes
perpendicular to the axis. Suppose that the planes are a distance 𝑕 apart and that the
circles that form the top and bottom of the frustum have radii 𝑎 and 𝑏. Then the volume
1
of the frustum equals 3 𝜋𝑕 𝑎2 + 𝑎𝑏 + 𝑏 2 . Let 𝑙 be the slant height of the frustum; that is,

the length of the part of a generator between the top and bottom of the frustum. Then
the area of the curved surface of the frustum equals 𝜋 (𝑎 + 𝑏)𝑙.

F-SIGMA: In a topological space 𝑿, a subset which is the countable union of closed sets.

FUBINI'S THEOREM: Suppose 𝑋 and 𝑌 are measure spaces, and suppose that 𝑋 × 𝑌 is
given the maximal product measure (which is the only product measure if 𝑋 and 𝑌 are
𝜍-finite). Fubini's theorem states that if 𝑓(𝑥, 𝑦) is 𝑋 × 𝑌 integrable, meaning that it is
measurable and

then

The first two integrals are iterated integrals with respect to two measures, respectively,
and the third is an integral with respect to the maximal product of these two measures.

FULLY-CLOSED MAPPING: A continuous mapping 𝒇: 𝑿 → 𝒀 with the following property:

For any point 𝒚 ∇ 𝒀 and for any finite family 𝑶𝟏 , 𝑶𝟐 , … … … , 𝑶𝒔 of open subsets of the
𝒔
space 𝑿 such that 𝒇−𝟏 𝒚 = 𝒊=𝟏 𝑶𝒊 , the set {𝒚} ∪ ⋃𝒔𝒊=𝟏 𝒇# 𝑶𝒊 is open.
Here 𝒇# 𝑶𝒊 denotes the small image of the set 𝑶𝒊 under the mapping 𝒇. Any fully-closed
mapping is closed. The inequality 𝐝𝐢𝐦 𝑿 ≤ 𝐦𝐚𝐱{𝐝𝐢𝐦 𝒀, 𝐝𝐢𝐦 𝒇} is valid for any fully-
closed mapping 𝒇: 𝑿 → 𝒀 of a normal space 𝑿. For this reason, fully-closed mappings can
be employed to isolate fairly wide classes of compacta with non-coinciding dimensions
𝒅𝒊𝒎 and 𝒊𝒏𝒅. Moreover, 𝐝𝐢𝐦 𝒀 ≤ 𝐝𝐢𝐦 𝑿 + 𝟏 irrespective of the multiplicity of the
mapping 𝒇. Let 𝒚 ∇ 𝒀, let 𝒇: 𝑿 → 𝒀 be a fully-closed mapping and let 𝑹(𝒇, 𝒀) be the
decomposition of 𝑿 the elements of which are all pre-images 𝒇−𝟏 (𝒚′ ) of the points, and
all points of 𝒇−𝟏 (𝒚). Then, for a regular space 𝑿, the quotient space 𝑿𝒀𝒇 of 𝑿 with respect
to the decomposition 𝑹(𝒇, 𝒀) is also regular; this property is characteristic of fully-
closed mappings in the class of closed mappings.
FUNCTION: A function is a rule that associates each member of one set with a unique
member of another set. The input number to the function is called the independent
variable, or argument. The set of all possible values for the independent variable is
called the domain. The output number is called the dependent variable. The set of all
possible values for the dependent variable is called the range. An important property of
functions is that for each value of the independent variable there is one and only one
value of the dependent variable. The notation 𝑓: 𝑆 → 𝑇, read as ‘𝑓 from 𝑆 to 𝑇′, is used.
If 𝑥 ∇ 𝑆, then 𝑓(𝑥) is the image of 𝑥 under 𝑓. The subset of 𝑇 consisting of those
elements that are images of elements of 𝑆 under 𝑓, that is, the set {𝑦 | 𝑦 = 𝑓(𝑥), for
some 𝑥 in 𝑆}, is the range of 𝑓. If 𝑓(𝑥) = 𝑦, it is said that 𝑓 maps 𝑥 to 𝑦, written
𝑓: 𝑥 → 𝑦. If the graph of 𝑓 is then taken to be 𝑦 = 𝑓(𝑥), it may be said that 𝑦 is a
function of 𝑥.
FUNCTIONAL: The functional Φ is an assignment of a scalar Φ (𝑓) to every vector 𝑓 in a
certain subset 𝐷Φ of the normed linear space 𝑁. 𝐷Φ is the domain of Φ.

FUNCTIONAL ANALYSIS: The part of modern mathematical analysis in which the basic
purpose is to study functions 𝒚 = 𝒇(𝒙) for which at least one of the
variables 𝒙 or 𝒚 varies over an infinite-dimensional space. In its most general form such
a study falls into three parts:
1) the introduction and study of infinite-dimensional spaces as such;
2) the study of the simplest functions, namely, when takes values in an infinite-
dimensional space and in a one-dimensional space (these are called functionals,
whence the name "functional analysis" ); and
3) the study of general functions of the type indicated — operators.
Linear functions 𝒙 ∇ 𝑿 ⇒ 𝒇 𝒙 = 𝒚 ∇ 𝒀 , i.e. linear operators, have been most
completely studied. Their theory is essentially a generalization of linear algebra to the
infinite-dimensional case. A combination of the approaches of classical analysis and
algebra is characteristic for the methods of functional analysis, and this leads to
relations between what are at first glance very distant branches of mathematics.
FUNCTIONAL EQUATION: An equation (linear or non-linear) in which the unknown is
an element of some specific (function) or abstract Banach space, that is, an equation of
the form
𝑃 𝑥 =𝑦

where 𝑃 𝑥 is some, generally speaking non-linear, operator transforming elements of

a 𝐵-space 𝑋 into elements of a space 𝑌 of the same type. If the functional equation
contains another numerical (or, in general, functional) parameter 𝜆, then instead of (1)
one writes

𝑃 𝑥, 𝜆 = 𝑦

where𝑥 ∇ 𝑋, 𝑦 ∇ 𝑌, 𝜆 ∇ Λ and Λ is the parameter space.

FUNCTIONAL OF A BOUNDARY VALUE PROBLEM: Consider the problem of finding a

curve through two points 𝑥1 , 𝑦1 and 𝑥2 , 𝑦2 whose length is minimum.
In general ,we wish to find a curve C joining the points 𝑥1 , 𝑦1 and 𝑥2 , 𝑦2 and having
equation 𝑦 = 𝑦 𝑥 ,where 𝑦 𝑥1 = 𝑦1 and 𝑦 𝑥2 = 𝑦2 such that for a given function
𝑓 𝑥, 𝑦, 𝑦′ .
𝑥2
∫𝑥 𝑓 𝑥, 𝑦, 𝑦′ 𝑑𝑥 is a stationary value or an extremum.
1

An integral such as above, which assume a definite value for functions of the type
𝑦 = 𝑦 𝑥 is called a functional.
FUNCTIONAL RELATION: A binary relation 𝑹 on a set 𝑨 satisfying 𝑹−𝟏 ∘ 𝑹 ⊆ 𝜟,
where 𝜟 is the diagonal of 𝑨. This means that (𝒂, 𝒃) ∇ 𝑹 and (𝒂, 𝒄) ∇ 𝑹 imply that 𝒃 = 𝒄,
that is, for each 𝒂 ∇ 𝑨 there is at most one 𝒃 ∇ 𝑨 such that (𝒂, 𝒃) ∇ 𝑹.
Thus, 𝑹 determines a function on 𝑨. When it satisfies 𝑹−𝟏 ∘ 𝑹 = 𝜟, this function is well-
defined everywhere and is one-to-one.
FUNCTION ELEMENT: An analytic function 𝑓 with domain 𝐷 is called a function element
and is denoted by 𝑓, 𝐷 .

FUNCTION FIELD ANALOGY: It was realized in the nineteenth century that the ring of
integers of a number field has analogies with the affine coordinate ring of an algebraic
curve or compact Riemann surface, with a point or more removed corresponding to the
'infinite places' of a number field. This thought is more precisely encoded in the theory
that global fields should all be treated on the same basis. The scheme goes further. Thus
elliptic surfaces over the complex numbers, also, have some quite strict analogies with
elliptic curves over number fields.
FUNCTION F SUMMABLE BY A MEASURE: A function 𝑓 is summable by a measure µ if
there is sequence (𝑓𝑛 ) ⊂ 𝑆(𝑋) such that

1. the sequence (𝑓𝑛 ) is a Cauchy sequence in 𝑆(𝑋);

𝑎. 𝑒
2. 𝑓𝑛 𝑓.
→

Note that

 Clear if a function is summable, then any equivalent function is summable as

well. Set of equivalent classes will be denoted by L1(X).
 If the measure µ is finite then any bounded measurable function is summable.
 This Lemma can be extended to the space of essentially bounded functions
𝐿∞ (𝑋), in other words 𝐿∞ (𝑋) ⊂ 𝐿1 (𝑋) for finite measures.
  If the measure µ is finite and 𝑓𝑛 ⇒ 𝑓 then 𝑑1 (𝑓𝑛 , 𝑓) → 0.
 Let (𝑓𝑛 ) and (𝑔𝑛 ) be Cauchy sequences in 𝑆(𝑋) with the same limit a.e., then
𝑑1 (𝑓𝑛 , 𝑔𝑛 ) → 0.

FUNCTIONS OF CLASS 𝒎: Let 𝐼 denote a real interval and let 𝑚 be a positive integer.
Then we say that a real-valued function 𝑓 defined on 𝐼is of class 𝑚 (or simple a Cm
function) if 𝑓 has a continuous derivative of 𝑚𝑡𝑕 order at every point 𝐼. In case 𝑓 is
differentiable an infinitely many number of times, it is said to be of class ∞ or a 𝐶 ∞ -
function. Also, the function 𝑓 is said to be analytic over 𝐼 if 𝑓 is single-valued and
possesses continuous derivatives of all orders at every point of 𝐼. An analytic function is
said to be of class ω or simply a 𝐶 ∞ - function.
FUNCTION THEORY: A branch of mathematics in which one studies the general
properties of functions. Function theory splits into two parts: the theory of functions of
a real variable and the theory of functions of a complex variable.
FUNDAMENTAL AXIOM OF ANALYSIS: If 𝑎1 , 𝑎2 , . ..is an increasing sequence in 𝑅 and
there exists an 𝐴 ∇ 𝑅 such that 𝑎𝑛 ≤ 𝐴 for all 𝑛, then there exists an 𝑎 ∇ 𝑅 such that
𝑎𝑛 → 𝑎 𝑎𝑠 𝑛 → ∞.
FUNDAMENTAL COUNTING PRINCIPLE: It is defined as the law for determining the
number of ways two or more operations can be performed together. If one operation
can be performed in 𝑚 ways and a second in 𝑛 ways, together they can be performed in
𝑚𝑛 ways.
FUNDAMENTAL MATRIX: The transition matrix 𝑋(𝑡) of solutions of a system of linear
ordinary differential equations
𝑥 = 𝐴 𝑡 𝑥; 𝑥 ∇ 𝑅 𝑛

normalized at the point 𝑡0 . The fundamental matrix is the unique continuous solution of
the matrix initial value problem

𝑋 = 𝐴 𝑡 𝑋; 𝑋 𝑡0 = 1

(𝐼 denotes the identity matrix) if the matrix-valued function 𝐴 𝑡 is locally summable

over some interval 𝐽 ⊂ 𝑅, 𝑡 ∇ 𝐽.

Every matrix 𝑀(𝑡) built of column-solutions 𝑥1 , 𝑥2 , … … … , 𝑥𝑚 of the system, where 𝑚 is

a natural number, is expressible as 𝑀 𝑡 = 𝑋 𝑡 𝑀(𝑡0 ) . In particular, every
solution 𝑥(𝑡) of system can be written in the form 𝑥 𝑡 = 𝑋 𝑡 𝑥0 .
FUNDAMENTAL PLANES: The three planes, osculating plane, normal plane and
rectifying plane associated with each point of a curve are called as fundamental planes.
These planes are mutually perpendicular and are determined by moving trihedral 𝑡, 𝑛, 𝑏
at the point.
The equations of fundamental planes are:
Osculating plane: it contains 𝑡 and 𝑛 is normal to 𝑏, its equation is 𝑅 − 𝑟 ∙ 𝑏 = 0.
Normal plane: It contains 𝑛 and 𝑏 and is normal to 𝑡, its equation is 𝑅 − 𝑟 ∙ 𝑡 = 0.
Rectifying plane: It contains 𝑏 and 𝑡 and is normal to 𝑛, its equation is 𝑅 − 𝑟 ∙ 𝑛 = 0.
FUNDAMENTAL SEQUENCE: A sequence 𝑓1 , 𝑓2 , … of vectors of a normed linear space 𝑁
is said to be fundamental if given any real number 𝜀 > 0, a natural number 𝑁 (𝜀) can be
indicated such that

𝑓𝑛 − 𝑓𝑛 +𝑚 < 𝜀, 𝑚 ≥ 0.

Whenever 𝑛 ≥ 𝑁 𝜀 .

FUNDAMENTAL SET OF SOLUTION OF THE EQUATION 𝐀𝐗 = 𝐎: Suppose the rank r of

the coefficient matrix 𝐴 is less than the number of unknowns 𝑛. In this case the given
equations have a set of 𝑛 − 𝑟 linearly independent solutions and every possible solution
is a linear combination of these 𝑛 − r solutions . This set of 𝑛 − r solutions is called a
fundamental set of solutions of the equation 𝐴𝑋 = 𝑂.

A set of linearly independent solutions 𝑥1 , 𝑥2 , … , 𝑥𝑘 of the system of homogeneous

equations 𝐴𝑋 = 𝑂 is called the fundamental system solution of 𝐴𝑋 = 𝑂, if every solution
𝑋 of 𝐴𝑋 = 𝑂 can be written as a linear combination of these vectors i.e., in the form
X=c1 x1 + c2 x2 + ⋯ + ck xk , where c1 , c2 , … , ck are suitable numbers.

FUNDAMENTAL THEOREM OF ALGEBRA: The fundamental theorem of algebra says that

an nth-degree polynomial equation has at least one root among the complex numbers. It
has exactly n roots when you include complex roots and you realize that a root may
occur more than once.
FUNDAMENTAL THEOREM OF ARITHMETIC: The fundamental theorem of arithmetic
says that any natural number can be expressed as a unique product of prime numbers.
FUNDAMENTAL THEOREM OF CALCULUS ALONG CURVES: Let 𝑈 ⊆ 𝐶 be a domain and
let 𝛾 ∶ [𝑎, 𝑏] → 𝑈 be a 𝐶1 curve. If 𝑓 ∇ 𝐶1 (𝑈), then
 𝑏
𝜕𝑓 𝑑𝛾1 𝜕𝑓 𝑑𝛾2
𝑓(𝛾(𝑏)) − 𝑓(𝛾(𝑎)) = (𝛾(𝑡)) · + (𝛾(𝑡)) · 𝑑𝑡.
𝜕𝑥 𝑑𝑡 𝜕𝑦 𝑑𝑡
𝑎

FUNDAMENTAL THEOREM OF INTEGRAL CALCULUS: Let 𝑓 be a continuous function on

[𝑎, 𝑏]. Let 𝜑 be a differentiable function on [𝑎, 𝑏] such that 𝜑 ′ 𝑥 = 𝑓 𝑥 ∀ 𝑥 ∇ 𝑎, 𝑏 .
Then
𝑏

𝑓 𝑥 𝑑𝑥 = 𝜑 𝑏 − 𝜑(𝑎)
𝑎

FUNDAMENTAL THEOREM OF GAME THEORY: Suppose that, in a matrix game, E(x, y) is

the expectation, where x and y are mixed strategies for the two players. Then
max min 𝐸(𝒙, 𝒚) = min max 𝐸(𝒙, 𝒚) .
𝑥 𝑦 𝑦 𝑥

This theorem is also known as the ‘Minimax Theorem’.

FUNDAMENTAL THEOREM ON HOMOMORPHISMS (FUNDAMENTAL HOMOMORPHISM

THEOREM): Given two groups 𝐺 and 𝐻 and a group homomorphism 𝑓 : 𝐺 → 𝐻, let 𝐾 be
a normal subgroup in 𝐺 and 𝜑 the natural surjective homomorphism 𝐺 → 𝐺/𝐾. If 𝐾 is a
subset of 𝑘𝑒𝑟(𝑓) then there exists a unique homomorphism 𝑕: 𝐺/𝐾 → 𝐻 such
that 𝑓 = 𝑕 𝜑.

The situation is described by the following commutative diagram:

By setting 𝐾 = 𝑘𝑒𝑟 𝑓 , we immediately get the first isomorphism theorem.

FUZZY COMPLEMENT: The fuzzy complement of A is a fuzzy set 𝐴 in 𝑈 whose

membership function is defined as
𝝁𝑨 𝑥 = 1 − 𝝁𝑨 𝑥 ∀ 𝑥 ∇ 𝐴.
FUZZY EQUAL SETS: We say fuzzy sets A and B are equal if and only if 𝝁𝑨 𝑥 =
𝝁𝑩 𝑥 ∀ 𝑥 ∇ 𝑈.
FUZZY NUMBER: A fuzzy number 𝑀 is a convex normalized fuzzy set ˜M of the real line R
such that
  it exists exactly one 𝑥0 ∇ 𝑅, 𝜇𝑀 (𝑥0 ) = 1 (𝑥0 is called the mean value of 𝑀 ).
 𝜇𝑀 (𝑥) is piecewise continuous.
In other words, A fuzzy number 𝐴 is a fuzzy set of the real line with a normal, (fuzzy)
convex and continuous membership function of bounded support. The family of fuzzy
numbers will be denoted by 𝐹.
FUZZY NUMBER OF LR-TYPE: A fuzzy number 𝑀 is of LR-type if there exist reference
functions L (for left) and R (for right), and scalars 𝛼 > 0, 𝛽 > 0, with
𝑚−𝑥
𝐿 ; 𝑖𝑓 𝑥 ≤ 𝑚
𝛼
𝜇𝑀 𝑥 = 𝑥−𝑚
𝑅 ; 𝑖𝑓 𝑥 ≥ 𝑚
𝛽
FUZZY POINT: Let 𝐴 be a fuzzy number. If 𝑠𝑢𝑝𝑝(𝐴) = {𝑥0 } then A is called a fuzzy point
and we use the notation 𝐴 = 𝑥0 .
FUZZY SET: A fuzzy set in a universe of discourse U is characterized by a membership
function 𝝁𝑨 (𝑥) that takes values in the interval 0, 1 . Therefore, a fuzzy set is a
generalization of a classical set by allowing the membership function to take any values
in the interval 0, 1 . In other words, the membership function of a classical set can only
take two values-zero and one, whereas the membership function of a fuzzy set is a
continuous function with range [0, 1]. A fuzzy set A in U may be represented as a set of
ordered pairs of a generic element x and its membership value, that is,
𝑨 = (𝑥, 𝝁𝑨 𝑥 ; 𝑥 ∇ 𝑼
FUZZY SUBSET: We say in fuzzy theory, 𝐵 contains 𝐴, denoted by 𝐴 ⊆ 𝐵, if and only if
𝝁𝑨 𝑥 ≤ 𝝁𝑩 𝑥 ∀ 𝑥 ∇ 𝑈.
G
Gal: It is the abbreviation for Galois group. (Also written as Γ .
GALOIS CORRESPONDENCE: Let 𝐾, 𝐿 and 𝑀 be fields satisfying 𝐾 ⊂ 𝑀 ⊂ 𝐿. Suppose
that 𝐿: 𝐾 is a Galois extension. Then so is 𝐿: 𝑀 . If in addition 𝑀: 𝐾 is normal, then
𝑀: 𝐾 is a Galois extension.
GALOIS, EVARISTE: Evariste Galois (1811 to 1832) was French mathematician who
made crucial contributions to group theory and applied this to the study of the
solvability of polynomial equations. He made major contributions to the theory of
equations before he died at the age of 20, shot in a duel. His work developed the
necessary group theory to deal with the question of whether an equation can be solved
algebraically.
GALOIS EXTENSIONOF A FIELD: A Galois extension is an algebraic extension of a
field that is normal and separable. The study of the automorphism group of such an
extension forms part of Galois theory.
GALOIS FIELD: A field with a finite number of elements. First considered by E. Galois.
GALOIS GROUP OF A FIELD EXTENSION: The Galois group 𝜞(𝑳: 𝑲) of a field extension
𝑳: 𝑲 is the group of all automorphisms of the field 𝑳 that fix all elements of the subfield
𝑲.
GALOIS GROUP OF A POLYNOMIAL: Let f be a polynomial with coefficients in some field
𝑲. The Galois group 𝜞𝑲(𝒇) of 𝒇 over 𝑲 is defined to be the Galois group 𝜞(𝑳: 𝑲) of the
extension 𝑳: 𝑲 , where 𝑳 is some splitting field for the polynomial 𝒇 over 𝑲.
GALOIS THEORY: In the most general sense, Galois theory is a theory dealing with
mathematical objects on the basis of their automorphism groups. For instance, Galois
theories of fields, rings, topological spaces, etc., are possible. In a narrower sense Galois
theory is the Galois theory of fields. The theory originated in the context of finding roots
of algebraic equations of high degrees. The familiar formula for solving equations of
degree two dates back to early Antiquity. Methods for solving cubic and quartic
equations were discovered in the 16th century. Unsuccessful attempts to find formulas
for solving quintic and higher-degree equations were made during the three centuries
which followed. It was finally proved by N.H. Abel in 1824 that there are no solutions in
radicals of the general equation of degree ≥ 𝟓. The next problem was to find necessary
and sufficient conditions to be satisfied by the coefficients of an equation for the latter
to be solvable in radicals, i.e. for it to be reducible to a chain of two-term equations of
the form 𝒙𝒏 − 𝒂 = 𝟎. This problem was solved by E. Galois; his results were exposed in
a letter on the eve of his death (1832), and published in 1846.
GAME THEORY: Game theory is the mathematical study of strategy games whose results
can be represented by a matrix showing the decisions of each player. Games may be
used to investigate problems in business, personal relationships, military manoeuvres
and other areas involving decision-making. One particular kind of game for which the
theory has been well developed is the matrix game.
GAMMA DISTRIBUTION: If 𝑥 is a random variable with p.d.f. given by
𝜆𝑣 𝑥 𝑣−1 𝑒 −𝜆𝑥
𝑓 𝑥 =
Γ(𝑣)
where 𝛤(𝑛) is a gamma function and 𝜆, 𝜈, 𝑥 > 0, then we say that X has a gamma
distribution with parameters 𝜆, 𝜈.
When 𝜈 = 1, 𝑥 𝜈−1 = 1, and 𝛤(𝜈) = 𝛤(1) = 1, so 𝑓(𝑥) reduces to 𝜆𝑒 −𝜆𝑥 which is the
exponential distribution.
∞
GAMMA FUNCTION: The function defined by 𝛤 𝑥 = ∫0 𝑡 𝑥 −1 𝑒 −𝑡 𝑑𝑡 for 𝑥 > 0.
Integration by parts yields 𝛤(𝑥 + 1) = 𝑥𝛤(𝑥), and 𝛤(1) = 1 so if 𝑛 is an integer
𝛤(𝑛) = (𝑛 − 1)!
GAUSS BONNET THEOREM: (DIFFERENTIAL GEOMETRY): Let a simply connected
region R be enclosed by a closed curve Γ composed of n smooth arcs
𝐴0 𝐴1 , 𝐴1 𝐴2, … , 𝐴𝑛−1 𝐴𝑛 𝐴𝑛 = 𝐴0 , making at the vertices exterior angles 𝛼1 , 𝛼2 , … , 𝛼𝑛 ;
then the excess of Γ defined as

𝑒𝑥. Γ = ∫2 𝐾𝑑𝑠 − 2𝜋 − 𝛼𝑟 − ∫𝑟 𝜅𝑔 𝑑𝑠,

𝑟=1

Where 𝜅 is some function of 𝑢, 𝑣 called the Gaussian Curvature and 𝜅𝑔 is the geodesic
curvature of the arcs.

GAUSS, CARL FRIEDRICH GAUSS: Carl Friedrich Gauss (1777 to 1855), perhaps the
greatest pure mathematician of all time, was a German mathematician who studied
errors of measurement (so the normal curve is sometimes called the Gaussian error
curve); developed a way to construct a 17-sided regular polygon with geometric
construction; developed a law that says the electric flux over a closed surface is
proportional to the charge inside the surface and studied the theory of complex
numbers. By the age of 24, he was ready to publish his Disquisitiones arithmeticae, a
book that was to have a profound influence on the theory of numbers. In this, he proved
the Fundamental Theorem of Arithmetic and the Fundamental Theorem of Algebra. In
later work, he developed the theory of curved surfaces using methods now known as
differential geometry. His work on complex functions was fundamental but, like his
discovery of non-Euclidean geometry, it was not published at the time. He introduced
what is now known to statisticians as the Gaussian distribution.
GAUSSIAN CURVATURE OF A SURFACE: The product of the principal curvatures of a
regular surface at a given point. If
𝒅𝒔𝟐 = 𝑬𝒅𝒖𝟐 + 𝟐𝑭𝒅𝒖𝒅𝒗 + 𝑮𝒅𝒗𝟐
is the first fundamental form of the surface and

𝑳𝒅𝒖𝟐 + 𝟐𝑴𝒅𝒖𝒅𝒗 + 𝑵𝒅𝒗𝟐

is the second fundamental form of the surface, then the Gaussian curvature can be
computed by the formula

𝐿𝑁 − 𝑀2
𝐾=
𝐸𝐺 − 𝐹 2
2
GAUSSIAN FUNCTION: The function 𝑓(𝑥) = 𝑒 –𝑥 which has the property which is the
function underlying the normal distribution.
GAUSSIAN INTEGER: A complex number whose real and imaginary parts are both
integers, so 𝑧 = 𝑎 + 𝑖𝑏 is known as a Gaussian integer if 𝑎, 𝑏 ∇ 𝑍.
GAUSSIAN RING: The ring of Gaussian integers or Gauss numbers, Z[i].
GAUSS HYPERGEOMETRIC EQUATION: The differential equation

d2 y dy
𝑥 1−𝑥 2
+ 𝑐− 𝑎+𝑏+1 𝑥 − aby = 0
dx dx

is known as Gauss’ hypergeometric equation or simple hypergeometric equation or

Gauss’a equation.

GAUSS-JORDAN ELIMINATION: Gauss-Jordan elimination is a method for solving a

system of linear equations. The method involves transforming the system so that the
last equation contains only one variable, the next-to- last equation contains only two
variables, and so on.. The result of this systematic method is that the augmented matrix
is transformed into reduced echelon form. As a method for solving simultaneous linear
equations, Gauss–Jordan elimination in fact requires more work than Gaussian
elimination followed by back-substitution, and so it is not in general recommended.

GAUSS–LUCAS THEOREM: If 𝑃 is a (non constant) polynomial with complex coefficients,

all zeros of 𝑃′ belong to the convex hull of the set of zeros of 𝑃.

GAUSS–MARKOV THEOREM: In a linear regression model in which the errors have zero
mean, are uncorrelated, and have equal variances the best linear unbiased estimators of
the coefficients are the least squares estimators. Here, ‘best’ means that it has minimum
variance amongst all linear unbiased estimators.
GAUSS MEAN VALUE THEOREM (COMPLEX ANALYSIS): If 𝑓(𝑧) is analytic in a domain 𝐷
and if the circular domain 𝑧 − 𝑧0 ≤ 𝜌 is contained in 𝐷, then

2𝜋
1
𝑓 𝑧0 = 𝑓 𝑧0 + 𝜌𝑒 𝑖𝜃 𝑑𝜃
2𝜋 0

That is to say, the value of 𝑓 𝑧 𝑎𝑡 𝑧0 is equal to the average of its value of the boundary
of the circle 𝑧 − 𝑧0 = 𝜌

GAUSS NUMBER: A complex integer 𝒂 + 𝒃𝒊, where 𝒂 and 𝒃 are arbitrary rational
integers. Geometrically, the Gauss numbers form the lattice of all points with integral
rational coordinates on the plane. Such numbers were first considered in 1832 by C.F.
Gauss in his work on biquadratic residues. He also discovered the properties of the
set 𝜞 of complex integers.
𝛤 is an integral domain; its units (i.e. divisors of the unit element) are 1, −1, 𝑖, −𝑖 and
there are no other units. One kind of primes (i.e. numbers that cannot be decomposed
into a non-trivial product) of 𝛤 (the Gaussian primes) are the numbers of the form
𝛼 = 𝑎 + 𝑏𝑖, the norms (moduli) 𝑁(𝛼) = 𝑎2 + 𝑏 2 = 𝑝 of which are rational prime
numbers 𝑝 of the form 4𝑛 + 1 or 𝑝 = 2; the other kind are rational prime numbers of
the form 4𝑛 + 3. Examples of Gaussian primes are 1 + 𝑖, 1 + 2𝑖, 3 + 4𝑖, 3, 7, etc.
Any number in 𝛤 can be uniquely decomposed into a product of primes in 𝛤, up to units
and ordering. Domains with this property are called unique factorization domains or
Gaussian rings.

In the theory of biquadratic residues the Gaussian numbers were the first simple and
important instance of an extension of the field of rational numbers.

GAUSS PRINCIPLE(PRINCIPLE OF LEAST FORCING): One of the fundamental and most

general differential variational principles of classical mechanics, established by C.F.
Gauss and expressing an extremum property of a real motion of a system in the class of
admissible motions, corresponding to the ideal constraints imposed on the system and
to the conditions of constancy of positions and velocities of the points in the system at a
given moment of time.

According to the Gauss principle, "the motion of a system of material points, constrained
in an arbitrary manner, and subjected to arbitrary forces at any moment of time, takes
place in a manner which is as similar as possible to the motion that would be performed
by these points if they were free, i.e. with least-possible forcing — the measure of
forcing during the time dt being defined as the sum of the products of the mass of each
point and the square of the distance of the point from the position which it would
occupy if it were free".

The Gauss principle is equivalent with the d'Alembert–Lagrange principle and is

applicable both to holonomic and to non-holonomic systems. It has been generalized in
various ways such as to systems subject to non-ideal constraints, as well as to the case
of continuous media.

GAUSS RECIPROCITY LAW: A relation connecting the values of the Legendre symbols
(𝒑/𝒒) and (𝒒/𝒑) for different odd prime numbers 𝒑 and 𝒒 . In addition to the principal
reciprocity law of Gauss for quadratic residues, which may be expressed as the relation
(𝑝𝑞)(𝑞𝑝) = (−1)(𝑝−1)/2⋅(𝑞−1)/2 ,

there are two more additions to this law, viz.:

2 −1)/8
(−1/𝑝) = (−1)(𝑝−1)/2 𝑎𝑛𝑑(2/𝑝) = (−1)(𝑝 .

The reciprocity law for quadratic residues was first stated in 1772 by L. Euler. A.
Legendre in 1785 formulated the law in modern form and proved a part of it. C.F. Gauss
in 1801 was the first to give a complete proof of the law ; he also gave no less than eight
different proofs of the reciprocity law, based on various principles, during his lifetime.

Attempts to establish the reciprocity law for cubic and biquadratic residues led Gauss to
introduce the ring of Gaussian integers.

GAUSS–SEIDEL ITERATIVE METHOD: A technique for solving a set of 𝑛 linear equations

in 𝑛 unknowns. If the system is summarized by 𝑨𝒙 = 𝒃, then taking initial values as
(1) 𝑏
𝑥𝑖 = 𝑎 𝑖 , it uses the iterative relation
𝑖𝑖

(𝑘) (𝑘−1)
(𝑘)
𝑏𝑖 − 𝑗 <𝑖 𝑎𝑖𝑗 𝑥𝑗 − 𝑗 >𝑖 𝑎𝑖𝑗 𝑥𝑗
𝑥𝑖 =
𝑎𝑖𝑗
so it uses the new values immediately they are available, unlike the Jacobi method.
GAUSS’S LEMMA: (a) Let p(x) be a polynomial with integer coefficients. Then if p(x) can
be factorized using rational numbers, p(x) can be factorized using only integers.
(b) Let 𝑔 and 𝑕 be polynomials with integer coefficients. If 𝑔 and 𝑕 are both primitive
then so is 𝑔𝑕.
GAUSS TEST: Let 𝑢𝑛 be a series of positive terms such that
𝑢𝑛 𝑎 𝑏
= 1+ +
𝑢𝑛 +1 𝑛 𝑜(𝑛2 )
Then
1. 𝑢𝑛 converges if 𝑎 > 1 or 𝑎 = 1 and 𝑏 > 1.
2. 𝑢𝑛 diverges if 𝑎 < 1 or 𝑎 = 1 and 𝑏 ≤ 1.
G-DELTA SET: In a topological space 𝑿 a subset which is the countable intersection of
open sets.
GELFAND–MAZUR LEMMA: If 𝑋 is a complex Banach division algebra, then 𝑋 is
isometrically isomorphic to 𝐶.
GENERAL ABELIAN GROUPS: An Abelian group is, in general, an extension of a torsion
group by a torsion-free group. A torsion group T is called bounded if there is an integer
𝑛 such that 𝑡 𝑛 = 1 for all 𝑡 ∇ 𝑇. Suppose there is a torsion group T. Then T is a direct
summand of an Abelian group G which contains T as its maximal torsion subgroup if
and only if T is the direct product of a divisible group and a bounded group.
GENERAL ASSOCIATIVE LAW: Let 𝑨,∗ be a semigroup, and let 𝒙, 𝒚, 𝒛 and w be elements
of A. We can use the associative property of ∗ to show that the value of a product
involving 𝒙, 𝒚, 𝒛, 𝒘 is independent of the manner in which that product is bracketed,
though it generally depends on the order in which 𝒙, 𝒚, 𝒛 and 𝒘 occur in that product
(unless that binary operation is also commutative). For example,
(𝒙 ∗ (𝒚 ∗ 𝒛)) ∗ 𝒘 = ((𝒙 ∗ 𝒚) ∗ 𝒛) ∗ 𝒘 = (𝒙 ∗ 𝒚) ∗ (𝒛 ∗ 𝒘)
= 𝒙 ∗ (𝒚 ∗ (𝒛 ∗ 𝒘)) = 𝒙 ∗ ((𝒚 ∗ 𝒛) ∗ 𝒘)
GENERALIZED ANALYTIC FUNCTION: A function 𝒘 𝒛 = 𝒖 𝒙, 𝒚 + 𝒊𝒗(𝒙, 𝒚) satisfying a
system
𝝏𝒖 𝝏𝒗 𝝏𝒚 𝝏𝒚
− + 𝒂𝒖 + 𝒃𝒗 = 𝟎, + + 𝒄𝒖 + 𝒅𝒗 = 𝟎
𝝏𝒙 𝝏𝒚 𝝏𝒙 𝝏𝒙

with real coefficients 𝑎, 𝑏, 𝑐, 𝑑 that are functions of the real variables 𝑥 and 𝑦.

GENERALIZED HAHN-BANACH THEOREM (FOR REAL VECTOR SPACES): Let 𝑞 be a

convex functional on a real linear space 𝑆, and Φ be a linear functional defined on a
subspace 𝑀 of 𝑆, such that

Φ 𝑓 ≤𝑞 𝑓 , for all 𝑓 𝜖 𝑀.
Then Φ can be extended to a linear functional Ψ defined on the whole of 𝑆, much that
Ψ 𝑕 ≤ 𝑞 𝑕 , for every 𝑕 ∇ 𝑀.

GENERALIZED NILPOTENT GROUP: A group in one of the generalized nilpotent classes

of groups. A class of groups is called generalized nilpotent if it contains all nilpotent
groups and if its intersection with the class of finite groups is the class of all finite
nilpotent groups. Quite a number of classes of generalized nilpotent groups have been
considered; principally, the connections between them have been studied. The most
important classes of generalized nilpotent groups are the class of locally nilpotent
groups, the classes of nil groups and groups with a normalizer condition. The majority
of classes of generalized nilpotent groups were introduced in studying various
properties of central or normal series and systems of subgroups.
GENERALIZED SOLVABLE GROUP: A group from one of the generalized solvable classes
of groups. A class of groups is called generalized solvable if it contains all solvable
groups and if its intersection with the class of finite groups is the class of all finite
solvable groups. Quite a number of classes of generalized solvable groups have been
considered; principally, the connections between them have been studied. The most
important class of generalized solvable groups is the class of locally solvable. Other
classes have been introduced in the study of various properties of normal and
subnormal series.
GENERAL LINEAR GROUP: The general linear group of degree 𝒏 is the group of
all (𝒏 × 𝒏) invertible matrices over an associative ring 𝑲 with a unit; the usual symbols
are 𝑮𝑳𝒏 (𝑲) 𝒐𝒓 𝑮𝑳(𝒏, 𝑲).
GENERAL SOLUTION OF A DIFFERENTIAL EQUATION: A function containing 𝑛 distinct
arbitrary constants which satisfies an 𝑛-th order differential equation is said to be a
general solution. It is obtained as the sum of the complementary function and a
particular integral.
GENERATING FUNCTION OF AN INFINITE SEQUENCE: The power series 𝐹(𝑥), where
𝐹(𝑥) = 𝑓0 + 𝑓1 𝑥 + 𝑓2 𝑥 2 + 𝑓3 𝑥 3 + …, is the generating function for the infinite
sequence ⌌𝑓𝑛 ⌍. Note that the generating function for this sequence is 1/(1 − 𝑥 − 𝑥 2 ). The
use of generating functions enables sequences to be handled concisely and algebraically.
A difference equation for a sequence can lead to an equation for the corresponding
generating function, and the use of partial fractions, for example, may then lead to a
formula for the n-th term of the sequence.
GENERATOR MATRIX (CODING THEORY): A generator matrix G for a linear code 𝐶 is a
matrix whose rows form a basis for 𝐶. A generator matrix is said to be in standard form
if it is of the form (𝐼𝑘 | 𝑋), where 𝐼𝑘 denotes the 𝑘 × 𝑘 identity matrix.
GENUS: Genus is the maximum number of times a surface can be cut along simple closed
curves without the surface separating into disconnected parts.
Genus of an entire function: The integer equal to the larger of the two
numbers 𝒑 and 𝒒 in the representation of the entire function 𝒇(𝒛) in the form
∞
𝝀 𝑸(𝒛)
𝒛 𝒛 𝒛𝟐 𝒛𝒑
𝒇 𝒛 =𝒛 𝒆 𝟏− 𝒆𝒙𝒑 + 𝟐 + ⋯+ 𝒑
𝜶𝒌 𝜶𝒌 𝟐𝜶𝒌 𝒑𝜶𝒌
𝒌=𝟏

where 𝑞 is the degree of the polynomial 𝑸(𝒛) and 𝑝 is the least integer satisfying the
condition
∞
1
𝑝+1
<∞
𝛼𝑘
𝑘=1

The number 𝑝 is called the genus of the product appearing in formula.

GEODESIC CIRCLE: The set of points on a metric two-dimensional manifold whose

distance from a fixed point 𝑶 is a constant 𝒓. A special case is a circle in the Euclidean
plane.

If 𝑟 is small, a geodesic circle on a regular surface and, in general, in a two-dimensional

Riemannian space is a simple closed curve (not necessarily of a constant geodesic
curvature); each one of its points may be connected with 𝑂 by a unique shortest line
(the radius or radial geodesic), forming a right angle with the geodesic circle; a geodesic
circle bounds a convex region. If 𝑟 → 0, the length 𝑙 of a geodesic circle is connected with
the Gaussian curvature 𝐾 at the point 𝑂 by the relation

2𝜋𝑟 − 𝑙 𝜋𝐾
→
𝑟3 3

If 𝑟 is large, more than one radial geodesic may lead to the same point on the geodesic
circle, the circle may bound a non-convex region and may consist of several
components. A geodesic circle is frequently employed in studies in global geometry. A
geodesic circle in the sense of Darboux is a closed curve of constant geodesic curvature.
It is a stationary curve for the isoperimetric problem. It coincides with an ordinary
geodesic circle on surfaces of constant curvature.

GEODESIC CO-ORDINATES: If the parametric curves are orthogonal and one of the
families of parametric curves is geodesic then the co-ordinates of the any point on the
surface are called a set of geodesic co-ordinates. Since one of the parametric curves is
assumed to be geodesic and other is arbitrary therefore the geodesic co-ordinate can be
introduced in an arbitrary number of ways.

GEODESIC CURVATURE AT A POINT OF A CURVE 𝜸 = 𝒓(𝒕) ON A SURFACE : The rate of

rotation of the tangent to 𝛾 around the normal 𝒏 to 𝐹, i.e. the projection on 𝒏 of the
vector of the angular rate of rotation of the tangent moving along 𝛾. It is assumed
that 𝛾 and 𝐹 are regular and oriented, and that the velocity is taken relative to the arc
length 𝑠 along 𝛾. The geodesic curvature can be defined as the curvature of the
projection of 𝛾 on the plane tangent to 𝐹 at the point under consideration. The geodesic
curvature is
𝒓′ , 𝒓′′ , 𝒏
𝒌𝒈 =
𝒓′ 𝟑

where a prime denotes differentiation with respect to 𝜏.

GEODESIC (DIFFERENTIAL GEOMETRY): A geodesic curve follows the shortest distance

between two points through a particular space. For example, in Euclidian space a
straight line is the geodesic between two points. Along the surface of the Earth, a great
circle route is the geodesic. Therefore, a curve on a surface, joining two given points,
that is the shortest curve between the two points is called a geodesic. On a sphere, a
geodesic is an arc of a great circle through the two given points. This arc is unique
unless the two points are antipodal.

GEODESIC FLOW (DIFFERENTIAL GEOMETRY): A Geodesic flow is a flow on a tangent

bundle TM of a manifold M, generated by a vector field whose trajectories are of the

form where is a geodesic.

GEODESIC LINE: The notion of a geodesic line is a geometric concept which is a

generalization of the concept of a straight line in Euclidean geometry to spaces of a
more general type. The definitions of geodesic lines in various spaces depend on the
particular structure (metric, line element, linear connection) on which the geometry of
the particular space is based. In the geometry of spaces in which the metric is
considered to be specified in advance, geodesic lines are defined as locally shortest. In
spaces with a connection, a geodesic line is defined as a curve for which the tangent
vector field is parallel along this curve. In Riemannian geometry, where a metric in the
tangent space at each point of the considered manifold is given, while the lengths of
lines are obtained by subsequent integration, geodesic lines are defined as extremals of
the length functional.
GEODESIC MANIFOLD AT A POINT x: A submanifold 𝑴𝒌 of a smooth
manifold 𝑴𝒏 (Riemannian or with an affine connection) such that the geodesic lines
of 𝑴𝒏 that are tangent to 𝑴𝒌 at x have a contact of at least the second order with 𝑴𝒌 .
This requirement is fulfilled at all points if any geodesic in 𝑴𝒌 is also a geodesic in 𝑴𝒏 .
Such geodesic manifolds 𝑴𝒌 are called totally geodesic manifolds.
GEODESIC MAPPING (DIFFERENTIAL GEOMETRY): A surface S is said to be mapped
geodesically onto a surface S* if there is a differentiable homemorphism of S onto S*
such that geodesic on S go over into geodesic on S*.

GEODESIC TANGENT (DIFFERENTIAL GEOMETRY): Let 𝑃 be a point on the curve 𝐶.

Then geodesic tangent of curve 𝐶 at point 𝑃 is the geodesic which touches the curve at
𝑃. Thus geodesic tangent at any point on a curve is the geodesic which touches the curve
at the point.

GEODESIC (TENSORS): We have ds 2 = g ij dx i dx j . The length of are from the point P0(t0)
to the point P1(t1) is given by

1/2
t dx i dx j
s = ∫t 1 g ij dt.
0 dt dt

Consider all the curves passing through two fixed point P0 and P1. If any of these curves
for which the distance P0P1 measured along the curve is stationary, then that curve is
called geodesic.

GEODESIC TRIANGLE: A figure consisting of three different points together with the
pairwise-connecting geodesic lines. The points are known as the vertices, while the
geodesic lines are known as the sides of the triangle. A geodesic triangle can be
considered in any space in which geodesics exist. If the sides of a geodesic triangle
situated in a region homeomorphic to an open disc constitute a simply-closed contour,
then the interior domain is added to the geodesic triangle. On a regular surface the sum
of the angles of a geodesic triangle minus π (the excess of the triangle) is equal to the
total curvature of the interior region.

Given a geodesic triangle in a metric space, a plane triangle with the same side lengths is
often considered. This makes it possible to introduce various concepts of an angle
between two shortest lines in metric spaces. In the two-dimensional case, after an angle
measurement has been introduced, it is possible to introduce the total curvature as a set
function expressed in terms of the excess of geodesic triangles. Nets of geodesic
triangles serve as a source of the approximation of metrics by the polyhedral metrics.

Estimates are available of the difference between the angle of a geodesic triangle in the
space under consideration and the respective angle in a triangle with the same side
lengths in a plane or on a surface of constant curvature

GEOMETRIC DISTRIBUTION: It is the discrete probability distribution for the number of

experiments required to achieve the first success in a sequence of independent
experiments, all with the same probability p of success. Consider a random experiment
where the probability of success on each trial is 𝑝. We will keep conducting the
experiment until you see the first success; let 𝑋 be the number of failures that occur
before the first success. Then 𝑋 is a discrete random variable with the geometric
distribution. Its probability function is:
𝑃 𝑋 = 𝑖 = 𝑝 (1 − 𝑝)𝑖
The expectation of 𝑋 is (1 − 𝑝)/𝑝, and the variance is (1 − 𝑝)/𝑝2 .
GEOMETRIC MEAN: The geometric mean of a group of 𝑛 numbers (𝑎1 , 𝑎2 , 𝑎3 , . . . 𝑎𝑛 ) is
equal to 𝑛 (𝑎1 𝑎2 𝑎3 . . . 𝑎𝑛 )
GEOMETRIC SEQUENCE: A geometric sequence is a sequence of numbers of the form
𝑎, 𝑎𝑟, 𝑎𝑟 2 , 𝑎𝑟 3 , . . . 𝑎𝑟 𝑛
The ratio between any two consecutive terms is a constant.
GEOMETRIC SERIES A geometric series is a sum of a geometric sequence:
𝑆 = 𝑎 + 𝑎𝑟 + 𝑎𝑟 2 + 𝑎𝑟 3 + . . . +𝑎𝑟 𝑛
In a geometric series the ratio of any two consecutive terms is a constant (𝑟). The sum
of the 𝑛 terms of the geometric series above is
 𝑎 1 − 𝑟𝑛
𝑆𝑛 = 𝑎 + 𝑎𝑟 + 𝑎𝑟 2 + 𝑎𝑟 3 + . . . +𝑎𝑟 𝑛−1 =
1−𝑟
GEOMETRY: Geometry is the area of mathematics related to the study of points and
figures, and their properties. Geometry is the study of shape and size. Euclidian
geometry has a rigorously developed logical structure. Three basic undefined terms are
point, line, and plane. A point is like a tiny dot: it has zero height, zero width, and zero
thickness. A line goes off straight in both directions. A plane is a flat surface, like a
tabletop, extending off to infinity. Euclid developed some basic postulates and then
proved theorems based on these. Examples of postulates used in modern versions of
Euclidian geometry are “Two distinct points are contained in one and only one line” and
“Three distinct points not on the same line are contained in one and only one plane.”
The geometry of flat figures is called plane geometry, because a flat figure is contained
in a plane. The geometry of figures in three dimensional space is called solid geometry.
Geometry of numbers

GEOMETRIC NUMBER THEORY: The branch of number theory that studies number-
theoretical problems by the use of geometric methods. Geometry of numbers in its
proper sense was formulated by H. Minkowski in 1896 in his fundamental monograph.
The starting point of this science, which subsequently became an independent branch of
number theory, is the fact that certain assertions which seem evident in the context of
figures in an 𝑛-dimensional Euclidean space have far-reaching consequences in number
theory.

GF: It is the abbreviation for Galois field.

GL: It is the abbreviation for general linear group.
glb: It is the abbreviation for greatest lower bound. (Also written as inf).
GÖDEL, KURT (1906–78): Kurt Gödel was born in Brno, was at the University of Vienna
from 1930 until he immigrated to the United States in 1940. Kurt Gödel was a
mathematician who showed that the consistency of elementary arithmetic could not be
proved from within the system itself. This result followed from his proof that any formal
axiomatic system contains undecidable propositions. It undermined the hopes of those
who had been attempting to determine axioms from which all mathematics could be
deduced.
GOLAY CODES: Let 𝐺 be a 12 × 24 matrix 𝐺 = (𝐼12 | 𝐴) where 𝐴 is a special 12 × 12
matrix containing only 0’s and 1’s. Then a linear binary code with generator matrix 𝐺 is
called an extended binary Golay code and is denoted as 𝐺24 . Let 𝐺 be the 12 × 23
matrix defined by 𝐺 = (𝐼12 | 𝐴) where 𝐴 is obtained from 𝐴 by deleting the last column
of 𝐴. The binary linear code with generator 𝐺 is called the binary Golay code and is
denoted as 𝐺23 .
GOLDBACH, CHRISTIAN (1690–1764): Christian Goldbach was a Mathematician born in
Prussia, who later became professor in St Petersburg. Goldbach’s conjecture, for which
he is remembered, was proposed in 1742 in a letter to Euler.
GOLDBACH CONJECTURE: The conjecture that every even integer greater than 2 is the
sum of two primes. Till now, it is neither proved nor disproved, Goldbach’s conjecture
remains one of the most famous unsolved problems in number theory.
GOLDBACH PROBLEM: One of the well-known problems in number theory: To give a
proof that every odd integer equal to or larger than 7 can be written as the sum of three
prime numbers. It was posed in 1742 by Goldbach in a communication to L. Euler. Euler
replied by pointing out that in order to solve this problem it is sufficient to prove that
every even number greater than 4 is the sum of two prime numbers. All attempts to
solve the problem remained unsuccessful for a long time. G.H. Hardy and J.E. Littlewood
in 1923 succeeded in showing that if certain theorems concerning Dirichlet L-
functions (which have not been proved till now) are valid, then any sufficiently large
odd number is the sum of three prime numbers. I.M. Vinogradov in 1937 discovered a
new method in analytic number theory — the method of estimating trigonometric sums
involving prime numbers — and applied this method to prove an asymptotic formula
for the number of representations of odd numbers as sums of three prime numbers.
This formula implies that each sufficiently large odd number is the sum of three prime
numbers. This is one of the major achievements of modern mathematics. Vinogradov's
method offers a way for solving several problems of a much more general nature.

GOLDEN RATIO: The division of a line segment 𝒂 into two parts the greater of which, 𝒙,
is the mean proportional between the whole segment a and the smaller part 𝒂 − 𝒙, i.e.
𝑎: 𝑥 = 𝑥: (𝑎 − 𝑥).

To find 𝑥, one has to solve a quadratic equation,

𝑥 2 + 𝑎𝑥 − 𝑎2 = 0,

the positive solution of which is

 −1 + 5𝑎
𝑥= ≈ 0.62𝑎.
2

Condition (1) may also be written as

𝑥 𝑥 𝑎
1+ = 1 𝑜𝑟 𝑥 = 𝑥
𝑎 𝑎 1+𝑎

or

1
𝑥=𝑎 ,
1
1+ 1
1+1+⋯

i.e. representing x as a continued fraction, the convergents of which are

1 1 2 3 5 8 13
, , , , ,
1 2 3 5 8 13 21
, , …,where 1,1,2,3,5,8,13,21, …, are the Fibonacci numbers.

GOLDEN–THOMPSON INEQUALITY: This inequality was proved independently

by Golden (1965) and Thompson (1965), says that for Hermitian matrices A and B,

where 𝑡𝑟 is the trace, and 𝑒 𝐴 is the matrix exponential.

GOLDSTINE THEOREM: Let 𝑋 be a Banach space, then the image of the closed unit
ball 𝐵 ⊂ 𝑋 under the canonical imbedding into the closed unit ball 𝐵′′ of the bidual
space 𝑋 ′′ is weakly-dense.
GOOGOL: A fanciful name for the number 10100, written in decimal notation as a 1
followed by 100 zeros.
GOSSET, WILLIAM SEALY (1876–1937): William Sealy Gosset was a British statistician
best known for his discovery of the t-distribution.
GOURSET LEMMA (COMPLEX ANALYSIS): Give 𝜀 > 0 it is possible to divide the region
inside 𝐶 into finite number of meshes either complete square 𝐶𝑛 or partial square 𝐷𝑛
such that within each mesh there exists a point 𝑧0 for which

𝑓 𝑧 −(𝑧0 )
− 𝑓 ′ (𝑧0 ) < 𝜀 ∀𝑧 in the mesh.
𝑧−𝑧0

GRADIENT: The gradient of a curve at a point 𝑃 may be defined as equal to the gradient
of the tangent to the curve at 𝑃. The gradient of a multivariable function is a vector
consisting of the partial derivatives of that function. If 𝑓 (𝑥, 𝑦, 𝑧) is a function of three
variables, then the gradient of 𝑓, written as ∆𝑓, is the vector
𝜕𝑓 𝜕𝑓 𝜕𝑓
∆𝑓 = , ,
𝜕𝑥 𝜕𝑦 𝜕𝑧
GRADIENT METHOD: A method for the minimization of a function of several variables. It
is based on the fact that each successive approximation of the function 𝐹 is obtained
from the preceding one by a shift in the direction of the gradient of the function:
𝑥 𝑛 +1 = 𝑥 𝑛 − 𝛿𝑛 𝑔𝑟𝑎𝑑 𝐹 𝑥 𝑛

The parameter 𝛿𝑛 can be obtained, e.g., from the condition of the magnitude
𝐹 𝑥 𝑛 − 𝛿𝑛 𝑔𝑟𝑎𝑑 𝐹 𝑥 𝑛 being minimal.

GRAM’S DETERMINANT: If 𝑔1 , 𝑔2 , … . . , 𝑔𝑛 are n vectors of a Hilbert space 𝐻, then their

Gram’s determinant ℸ (𝑔1 , 𝑔2 , … , 𝑔𝑛 ), is

𝑔1 , 𝑔1 𝑔2 , 𝑔1 ⋯ 𝑔𝑛 , 𝑔1
𝑔1 , 𝑔2 𝑔2 , 𝑔2 ⋯ 𝑔𝑛 , 𝑔2
⋮ ⋮ ⋮ ⋮
𝑔1 , 𝑔𝑛 𝑔2 , 𝑔𝑛 ⋯ 𝑔𝑛 , 𝑔𝑛

GRAM’S DETERMINANT PRINCIPLE: Vectors 𝑔1 , 𝑔2,…, 𝑔𝑛 of a Hillbert space 𝐻 are linearly

dependent, if and only if

ℸ 𝑔1 , 𝑔2 , … , 𝑔𝑛 = 0.

They are linearly independent if and only if

ℸ 𝑔1 , 𝑔2 , … , 𝑔𝑛 ≠ 0,

and, in this case

ℸ 𝑔1 , 𝑔2 , … , 𝑔𝑛 > 0.

GRAM – SCHMIDT ORTHONORMALISATION PROCESS: If 𝑓1 , 𝑓2 , …. is a (possibly finite)

sequence of linearly independent vectors in a linear inner product space, than an
equivalent orthonormal sequence 𝑒1 , 𝑒2 , 𝑒3 , …, can be formed.

GRAM–SCHMIDT THEOREM: Let (𝑥𝑖 ) be a sequence of linearly independent vectors in an

inner product space 𝑉. Then there exists orthonormal sequence ⌌𝑒𝑛 ⌍1∞ such that
𝐿{𝑥1 , 𝑥2 , … , 𝑥𝑛 } = 𝐿{𝑒1 , 𝑒2 , … , 𝑒𝑛 } for all n. A separable Hilbert space can be identified
with either 𝑙2𝑛 𝑜𝑟 𝑙2 , in other words it has an orthonormal basis ⌌𝑒𝑛 ⌍ (finite or infinite)
such that
∞ ⌌𝑥, 2 ∞ ⌌𝑥, 𝑒𝑛 ⌍ 2 .
𝑥= 𝑛=1 𝑒𝑛 ⌍𝑒𝑛 𝑎𝑛𝑑 𝑥 = 𝑛=1

GRAPH: A number of vertices (nodes), some of which are joined by edges. The edge
joining the vertex 𝑈 and the vertex 𝑉 may be denoted by (𝑈, 𝑉) or (𝑉, 𝑈). The vertex-set,
that is, the set of vertices, of a graph 𝐺 may be denoted by 𝑉(𝐺) and the edge-set by
𝐸(𝐺). For example, the graph shown here on the left has 𝑉(𝐺) = {𝑈, 𝑉, 𝑊, 𝑋} and
𝐸(𝐺) = {(𝑈, 𝑉), (𝑈, 𝑊), (𝑉, 𝑊), (𝑊, 𝑋)}. In general, a graph may have more than one
edge joining a pair of vertices; when this occurs, these edges are called multiple edges.
Also, a graph may have loops—a loop is an edge that joins a vertex to itself.
GRAPH HOMEOMORPHISM: An equivalence relation on the set of graphs, characterizing
their geometric properties. The notion of a graph homeomorphism is defined as follows.
Subdivision of an edge (𝑎, 𝑏) of a graph 𝐺 is an operation involving the addition of a new
vertex 𝑣, the removal of (𝑎, 𝑏), and the addition of two new edges (𝑎, 𝑣) and (𝑣, 𝑏).
Geometrically, this operation consists in addition of some (interior) point 𝑣 on the
line (𝑎, 𝑏); this point then becomes a new vertex. A graph 𝐺′ is called a subdivision of a
graph 𝐺 if it can be obtained from 𝐺 by repeating the operation of edge subdivision
several times. Two graphs 𝐺1 and 𝐺2 are said to be homeomorphic if they have
isomorphic subdivisions.
GRAPHICAL SOLUTION OF AN LPP: The solution of a LPP obtained by graphical method
is called the graphical solution of LPP.

GRAPH ISOMORPHISM: An equivalence relation on the set of graphs. An isomorphic

mapping of a non-oriented graph to another one is a one-to-one mapping of the vertices
and the edges of one graph onto the vertices and the edges, respectively, of the other,
the incidence relation being preserved. Two graphs are said to be isomorphic if there
exists an isomorphic mapping of one of these graphs to the other. Isomorphic graphs
are usually not distinguished from one another. The number of pairwise non-
isomorphic graphs with a given number of vertices and a given number of edges is
finite. Isomorphism of oriented graphs, hypergraphs and networks can be defined in a
similar manner.
GRAPH OF A FUNCTION: For a real function 𝑓, the graph of 𝑓 is the set of all pairs (𝑥, 𝑦)
in 𝑹 × 𝑹 such that 𝑦 = 𝑓(𝑥) and 𝑥 is in the domain of the function. For many real
functions of interest, this gives a set of points that form a curve of some sort, possibly in
a number of parts that can be drawn in the plane. Such a curve defined by 𝑦 = 𝑓(𝑥) is
also called the graph of 𝑓 .
GRAPH OF A RELATION: Let 𝑅 be a binary relation on a set 𝑆, so that, when 𝑎 is related
to 𝑏, this is written 𝑎 𝑅 𝑏. The graph of 𝑅 is the corresponding subset of the Cartesian
product 𝑆 × 𝑆, namely, the set of all pairs (𝑎, 𝑏) such that 𝑎 𝑅 𝑏.
GRAPH THEORY: A branch of discrete mathematics, distinguished by its geometric
approach to the study of various objects. The principal object of the theory is a graph
and its generalizations. The first problems in the theory of graphs were solutions of
mathematical puzzles (the problem of the bridges of Königsberg, the disposition of
queens on a chessboard, transportation problems, the travelling-salesman problem,
etc.). One of the first results in graph theory was a criterion on the possibility of
traversing all edges of a graph without passing through any edge more than once; it was
obtained by L. Euler in 1736 in solving the problem of the bridges of Königsberg.
The four-colour problem, formulated in the mid-19th century, though a mere
amusement puzzle at first sight, led to studies of graphs of both theoretical and applied
interest. Certain studies in the mid-19th century contain results of importance to graph
theory, obtained by solving practical problems. In the 20th century, problems involving
graphs began to arise not only in physics, chemistry, electrical engineering, biology,
economics, sociology, etc., but also in mathematics itself — in topology, algebra,
probability theory, and number theory. At the beginning of the 20th century, graphs
were used to represent certain mathematical objects and to formally state various
discrete problems; besides the term "graph" , other terms such as map, complex,
diagram, network, labyrinth, were also employed. The first results, concerning
connectivity properties, planarity, and graph symmetry, which paved the way for a
number of novel directions of study in graph theory, appeared in the 1920s and 1930s.
The scope of research in graph theory was considerably extended in the late 1940s and
early 1950s, mainly as a result of the development of cybernetics and calculation
techniques. Interest in graph theory increased, and the range of problems dealt with by
the theory was considerably extended. It was shown, for individual classes of graphs
(trees, planar graphs, etc.), which had been studied before, that the solution of certain
problems was simpler for such graphs than for arbitrary graphs (finding conditions for
the existence of graphs with certain properties, establishment of graph isomorphism,
etc.).
GREAT CIRCLE: A great circle is a circle that is formed by the intersection of a sphere
and a plane passing through the center. A great circle is the largest circle that can be
drawn on a given sphere, and the shortest path along the sphere between two points is
a great circle. This great circle is unique unless the two points are antipodal.
GREATEST COMMON DIVISOR: The greatest common divisor of two natural numbers a
and b is the largest natural number that divides both a and b evenly (that is, with no
remainder) and is usually denoted as (𝑎, 𝑏). The greatest common divisor of 𝑎 and 𝑏 has
the property of being divisible by every other common divisor of 𝑎 and 𝑏. It is an
important theorem that there are integers 𝑠 and 𝑡 such that the greatest common
divisor can be expressed as 𝑠𝑎 + 𝑡𝑏. If the prime decompositions of 𝑎 and 𝑏 are known,
the greatest common divisor is easily found.
GREAT PICARD'S THEOREM: If an analytic function 𝑓 has an essential singularity at a
point 𝑤, then on any punctured neighborhood of 𝑤, 𝑓(𝑧) takes on all possible complex
values, with at most a single exception, infinitely often.
GREAT MEROMORPHIC PICARD'S THEOREM: If 𝑀 is a Riemann surface, 𝑤 a point
on 𝑀, 𝑷𝟏 (𝑪) = 𝑪 ∪ {∞} denotes the Riemann sphere and 𝑓 : 𝑀\{𝑤} → 𝑷𝟏 (𝑪) is a
holomorphic function with essential singularity at 𝑤, then on any open subset
of 𝑀 containing 𝑤, the function 𝑓(𝑧) attains all but at most two points of 𝑷𝟏 (𝑪) infinitely
often.
GREEDY ALGORITHM FOR THE MINIMAL SPANNING TREE PROBLEM:
 Select any node randomly, then connect it to the nearby distinct node.
 Recognize the unconnected node that is nearby to a connected node, put in this
arc to the network.
 Ties can be broken randomly. Such ties may point towards multiple optimal
solutions.
GREEN, GEORGE (1793–1841): George Green was a British mathematician who
developed the mathematical theory of electricity and magnetism. He had worked as a
baker, and was self-taught in mathematics; he published other notable mathematical
papers before beginning to study for a degree at Cambridge at the age of 40.
GREEN’S THEOREM:
 1. Let 𝒇 𝑥, 𝑦 = [u(𝑥, 𝑦), 𝑣(𝑥, 𝑦)] be a two-dimensional vector field, and let 𝐶 be a
closed path in the 𝑥, 𝑦 plane. Green’s theorem states that the line integral of 𝑓
around this path is equal to the following integral over the interior of the path 𝐶:
𝜕𝑢 𝜕𝑣
𝒇 𝑥, 𝑦 𝑑𝐶 = − 𝑑𝑥𝑑𝑦
𝜕𝑥 𝜕𝑦
𝐶 𝐼𝑛𝑡𝑒𝑟𝑖𝑜𝑟 𝑜𝑓 𝐶
2. If ∅, ∅′ are both single-valued and continuously differentiable scalar point
function such that ∆∅ and ∆∅’ are also continuously differentiable,
Then
∂∅
∆∅ . ∆∅′ dv = − ∅∆2 ∅′ dv − ∅ dS
V V S ∂n
∂∅
=− ∅′ ∆2 ∅ dv − ∅′ dS
V S ∂n
where S is closed surface bounding any simple- connected region, δn is an
element of the normal at any point on the boundary drawn into the region
considered, and V is the volume enclosed by 𝑆.

GREGORY, JAMES (1638–75): James Gregory was a Scottish mathematician who studied
in Italy. He obtained infinite series for certain trigonometric functions and was one of
the first to appreciate the difference between convergent and divergent series. He died
at the age of 36.

GROMOV-HAUSDORFF CONVERGENCE: Geodesic metric space is a metric space where

any two points are the endpoints of a minimizing geodesic.
GRONWALL'S INEQUALITY : Gronwall's inequality also called Grönwall's lemma
or Gronwall–Bellman inequality allows one to bound a function that is known to satisfy
a certain differential or integral inequality by the solution of the corresponding
differential or integral equation.

Let 𝐼 denote an interval of the real line of the form [𝑎,  ∞) or [𝑎, 𝑏] or [𝑎, 𝑏)
with 𝑎 < 𝑏. Let 𝛽 and 𝑢 be real-valued continuous functions defined on 𝐼.
If 𝑢 is differentiable in the interior 𝐼 and satisfies the differential inequality
then 𝑢 is bounded by the solution of the corresponding differential equation
𝑦 ′(𝑡) = 𝛽(𝑡) 𝑦(𝑡):

for all t ∇ I.

GRONWALL LEMMA: The Gronwall lemma is a fundamental estimate for (nonnegative)

functions on one real variable satisfying a certain differential inequality. The lemma is
extensively used in several areas of mathematics where evolution problems are studied
(e.g. partial and ordinary differential equations, continuous dynamical systems) to
bound quantities which depend on time. The most elementary version of the inequality
is stated in the following:
Let 𝜙: [0, 𝑇] → 𝑅 be a nonnegative differentiable function for which there exists a
constant 𝐶 such that 𝜙′(𝑡) ≤ 𝐶𝜙(𝑡)for all 𝑡 ∇ [0, 𝑇].
Then 𝜙(𝑡) ≤ 𝑒𝐶𝑡𝜙(0)for all 𝑡 ∇ [0, 𝑇].

GROUP: A group is a set 𝐺 of elements for which an operation ⊚ is defined that meets
the following properties:
(1) 𝑎, 𝑏 ∇ 𝐺 ⇒ 𝑎 ⊚ 𝑏 ∇ 𝐺. (Closure Axiom)
(2) The associative property holds:
𝑎 ⊚ 𝑏 ⊚ 𝑐 = 𝑎 ⊚ 𝑏 ⊚ 𝑐 ∀ 𝑎, 𝑏, 𝑐 ∇ 𝐺
(3) There is an identity element 𝑒 ∇ 𝐺 such that
𝑎⊚𝑒 =𝑎 = 𝑒⊚𝑎∀𝑎 ∇𝐺
(4) Each element 𝑎 ∇ 𝐺 has an inverse 𝑎 −1 ∇ 𝐺 such that
𝑎 ⊚ 𝑎−1 = 𝑒 = 𝑎−1 ⊚ 𝑎
If ⊚ is also commutative (that is, 𝑎 ⊚ 𝑏 = 𝑏 ⊚ 𝑎 ∀ 𝑎, 𝑏 ∇ 𝐺), then the group is called an
Abelian group.
For example, the real numbers form an Abelian group with respect to addition, and the
nonzero real numbers form an Abelian group with respect to multiplication.
The theory of groups can be applied to many sets other than numbers, and to operations
other than conventional multiplication.
GROUP ACTION OF A GROUP G ON A SET X: A map from 𝑿 × 𝑮 → 𝑿, written (𝒙, 𝒈) or 𝒙𝒈
satisfying
 (𝑥, 1𝐺 ) = 𝑥

(𝑥, 𝑔𝑕) = ((𝑥, 𝑔), 𝑕) .

For given 𝑔, the map 𝜌𝑔 : 𝑥 ↦ (𝑥, 𝑔) is a permutation of 𝑋, the inverse

mapping being 𝜌𝑔 −1 . The map 𝑔 ↦ 𝜌𝑔 is a homomorphism 𝜌: 𝐺 → 𝑆𝑋 where 𝑆𝑋 is the

symmetric group on 𝑋: conversely, every such homomorphism gives rise to an
action (𝑥, 𝑔) ↦ (𝑥)𝜌𝑔 . If the homomorphism 𝜌 is injective the action is faithful: 𝐺 may be
regarded as a subgroup of 𝑆𝑋 . In any case, the image of ρ is a permutation group on X.

If 𝑥 ∇ 𝑋, the orbit of 𝑥 is the set of points {(𝑥, 𝑔): 𝑔 ∇ 𝐺}. An action

is transitive if 𝑋 consists of a single orbit. An action is 𝑘 −fold transitive if for any
two 𝑘 −tuples of distinct elements (𝑥1 , … , 𝑥𝑘 ) and (𝑦1 , … , 𝑦𝑘 ) there is 𝑔 ∇ 𝐺 such
that 𝑦𝑖 = (𝑥𝑖 , 𝑔), 𝑖 = 1, … , 𝑘. An action is primitive if there is no non-trivial partition
of 𝑋 preserved by 𝐺. A doubly transitive action is primitive, and a primitive action is
transitive, but neither converse holds. 𝐹𝑜𝑟 𝑥 ∇ 𝑋, the stabiliser of 𝑥 is the subgroup 𝐺𝑥 =
{𝑔 ∇ 𝐺: (𝑥, 𝑔) = 𝑥}.
Burnside's Lemma states that the number 𝑘 of orbits is the average number of fixed
points of elements of 𝐺, that is, 𝑘 = |𝐺|−1 𝑔∇𝐺 |𝐹𝑖𝑥(𝑔)|, where 𝐹𝑖𝑥(𝑔) = {𝑥 ∇ 𝑋: 𝑥𝑔 =
𝑥} and the sum is over all 𝑔 ∇ 𝐺.
GROUPED DATA: Grouped data is a set of data is said to be grouped when certain groups
or categories are defined and the observations in each group are counted to give the
frequencies. For numerical data, groups are often defined by means of class intervals.
GROUPOID: A universal algebra with one binary operation. An important concept in the
theory of groupoids is that of isotopy of operations. On a set 𝑮, let there be defined two
binary operations, denoted by (⋅) and (∘); they are isotopic if there exist three one-to-
one mappings 𝜶, 𝜷 and 𝜸 of 𝑮 onto itself such that 𝒂 ⋅ 𝒃 = 𝜸 − 𝟏(𝜶𝒂 ∘ 𝜷𝒃) for
all 𝒂, 𝒃 ∇ 𝑮. A groupoid that is isotopic to a quasi-group is itself a quasi-group; a
groupoid with a unit element that is isotopic to a group, is also isomorphic to this group.
For this reason, in group theory the concept of isotopy is not used: For groups isotopy
and isomorphism coincide.
A groupoid with cancellation is a groupoid in which either of the equations 𝑎𝑏 =
𝑎𝑐, 𝑏𝑎 = 𝑐𝑎 implies 𝑏 = 𝑐, where 𝑎, 𝑏 and 𝑐 are elements of the groupoid. Any groupoid
with cancellation is imbeddable into a quasi-group. A homomorphic image of a quasi-
group is a groupoid with division, that is, a groupoid in which the equations 𝑎𝑥 =
𝑏 and 𝑦𝑎 = 𝑏 are solvable.
GROUP WITH UNIQUE DIVISION (R-GROUP): A group in which the equality 𝒙𝒏 =
𝒚𝒏 implies 𝒙 = 𝒚, where 𝒙, 𝒚 are any elements in the group and 𝒏 is any natural number.
A group 𝑮 is an R-group if and only if it is torsion-free and is such that 𝒙𝒏 𝒚 =
𝒚𝒙𝒏 implies 𝒙𝒚 = 𝒚𝒙 for any 𝒙, 𝒚 ∇ 𝑮 and any natural number 𝒏. An R-group splits into
the set-theoretic union of Abelian groups of rank 1 intersecting at the unit element. A
group is an R-group if and only if it is torsion-free and if its quotient group by the centre
is an R-group. Subgroups of an R-group, as well as direct and complete direct products
of R-groups, are R-groups. The following local theorem is valid for the class of R-groups:
If all finitely-generated subgroups of a group 𝑮 are R-groups, then 𝑮 itself is an R-group.
Free groups, free solvable groups and torsion-free locally nilpotent groups are R-
groups.

GRUNSKY'S THEOREM: Let 𝑓 be a univalent holomorphic function on the unit

disc 𝐷 such that 𝑓(0) = 0. Then for all 𝑟 ≤ 𝑡𝑎𝑛𝑕 𝜋/4, the image of the disc |𝑧| <
𝑟 is starlike with respect to 0, , i.e. it is invariant under multiplication by real numbers
in (0,1).
H
HADAMARD CODE: Let 𝐻𝑛 be a Hadamard matrix. A Hadamard code of length 𝑛, denoted
𝐻𝑎𝑑𝑛 , is the binary code derived from 𝐻𝑛 by replacing all −1 values with 0 in 𝐻𝑛 and
then taking all the rows of 𝐻𝑛 and their complements as codewords.
HADMARD FACTORIZATION THEOREM (COMPLEX ANALYSIS): If 𝑓(𝑧) is an entire
function of infinite order 𝜌, then

𝑓 𝑧 = 𝑧 𝑚 𝑒 𝑔(𝑧) 𝑃 (𝑧)

Where 𝑚 is the order of the (possible) zeroes of 𝑎𝑡 𝑧 = 0, 𝑔 is a polynomial of degree

not exceeding 𝑝 𝑎𝑛𝑑 𝑃(𝑧) is the canonical product associated with the sequence of non-
zeros of 𝑓 𝑧 .

HADAMARD MATRIX: A square matrix 𝐻𝑛 of dimension 𝑛 × 𝑛 is called a Hadamard

matrix if all of its entries are in {−1, 1} and it holds that 𝐻𝑛 · 𝐻𝑛𝑇 = 𝑛𝐼𝑛 .
HADAMARD MULTIPLICATION: The Hadamard product, or Schur product, of two 𝒎 ×
𝒏 matrices 𝑨 and 𝑩 is the 𝒎 × 𝒏 matrix 𝑨𝑩 with

(𝐴 ∘ 𝐵)𝑖𝑗 = 𝑎𝑖𝑗 𝑏𝑖𝑗 , 𝑖 = 1, … , 𝑚; 𝑗 = 1, … , 𝑛.

HADAMARD SPACE: Hadamard space is a complete simply connected space with

nonpositive curvature.

HADAMARD THEOREM: The function 𝜁(𝑠) has infinitely many zeros in the critical strip.

HADAMARD THREE CIRCLES PRINCIPLE (COMPLEX ANALYSIS): Suppose 𝑓(𝑧) is

analytic in the closed ring 𝑟1 ≤ 𝑧 ≤ 𝑟3.

Let 𝑟1 < 𝑟2 < 𝑟3 and 𝑀𝑖 be the maximum value of 𝑓(𝑧) on the circles 𝑧 = 𝑟𝑖 , 𝑖 =
1,2,3 . Then

𝑟 𝑟 𝑟
log ( 3 ) log ( 3 ) log ( 2 )
𝑟1 𝑟2 𝑟1
𝑀2 ≤ 𝑀1 . 𝑀3
The three circles theorem as a convexity theorem: Suppose 𝑓(𝑧) is analytic in the closed
ring 𝑟1 ≤ 𝑧 ≤ 𝑟2 . Let 𝑟1 < 𝑟2 < 𝑟3 and 𝑀 (𝑟𝑖 ) be the maximum value of 𝑓(𝑧) on the
circles 𝑧 = 𝑟𝑖 , 𝑖 = 1,2,3 . Then

log 𝑟3 − log 𝑟2 log 𝑟2 − log 𝑟1

log 𝑀 𝑟2 ≤ log 𝑀 𝑟1 + log 𝑀 𝑟3
log 𝑟3 − log 𝑟1 log 𝑟3 − log 𝑟1

HAHN–BANACH THEOREM: Let 𝐸 be a normed vector space, and let 𝐹 ⊆ 𝐸 be a

subspace. Let 𝜑 ∇ 𝐹 ∗ . Then there exists 𝜓 ∇ 𝐸 ∗ with ||𝜓|| ≤ ||𝜑|| and 𝜓(𝑥) = 𝜑(𝑥) for
each 𝑥 ∇ 𝐹.

Let 𝐸 be a normed vector space, and let 𝐹 be a subspace of 𝐸. For 𝑥 ∇ 𝐸, the following
are equivalent:

1. 𝑥 ∇ 𝐹 the closure of 𝐹;

2. for each 𝜑 ∇ 𝐸 ∗ with 𝜑(𝑦) = 0 for each 𝑦 ∇ 𝐹, we have that 𝜑(𝑥) = 0.

In simple words, let 𝑀 be a linear subspace of a normed linear space 𝑁 and let Φ be a
functional defined on 𝑀. Then Φ can be extended to a functional Φ0 defined on the
whole space 𝑁 such that Φ0 = Φ .

HAHN DECOMPOSITION: Let 𝜈 be a charge. There exist 𝐴, 𝐵 ∇ 𝐿, called a Hahn

decomposition of (𝑋, 𝜈), with 𝐴 ∩ 𝐵 = ∅, 𝐴 ∪ 𝐵 = 𝑋 and s.t. for any 𝐸 ∇ 𝐿, 𝜈 (𝐴⋂ 𝐸) ≥
0, 𝜈(𝐵⋂ 𝐸) ≤ 0. This need not be unique.

HAHN DECOMPOSITION THEOREM: Given a measurable space (𝑋, 𝛴) and a signed

measure 𝜇 defined on the ς-algebra 𝛴, there exist two measurable sets 𝑃 and 𝑁 in 𝛴
such that:

1. 𝑃 ∪ 𝑁 = 𝑋 and 𝑃 ∩ 𝑁 = ∅.

2. For each 𝐸 in 𝛴 such that 𝐸 ⊆ 𝑃 one has 𝜇(𝐸) ≥ 0; that is, 𝑃 is a positive
set for 𝜇.

3. For each 𝐸 in 𝛴 such that 𝐸 ⊆ 𝑁 one has 𝜇(𝐸) ≤ 0; that is, 𝑁 is a negative set
for 𝜇.
HALF-LIFE: The time it takes for a radioactive substance to decay to one-half its original
amount.
HALF PLANE: A half plane is the set of all points in a plane that lie on one side of an axis.
HALF-RANGE EXPANSIONS: In many physical problems, the function 𝑓(𝑥) is only known
over a finite interval, say 0 < 𝑥 < 𝐿. To express 𝑓(𝑥) as a Fourier series means that we
need to extend the function to be valid over all 𝑥. We can also extend the function in an
even manner to get a cosine series or as an odd function to get a sine series.

Figures:

(a) A function defined on a interval 0 < 𝑥 < 𝐿, where 𝐿 = 1.

(b) The even extension of the function onto the interval 0 < 𝑥 < 𝐿 is shown as a solid
curve and the periodic extension of period 2𝐿 is shown as a dot-dash curve.

(c) The odd extension of the function, solid curve, and the periodic extension of period
2𝐿, dot-dash curve.

The odd extension of 𝑓(𝑥) will generate a Fourier series that only involves sine terms
and is called the sine half-range expansion. It is given by

∞ 𝒏𝝅
𝒇 𝒙 = 𝒏=𝟏 𝒃𝒏 𝒔𝒊𝒏 𝑳 𝒙

2 𝐿 𝒏𝝅
where 𝑏𝑛 = 𝐿 ∫0 𝑓 𝑥 𝑠𝑖𝑛 𝒙𝑑𝑥 and 𝑛 is a positive integer.
𝑳
On the other hand, the even extension generates the cosine half-range expansion

∞ 𝒏𝝅
𝒇 𝒙 = 𝒂𝟎 + 𝒏=𝟏 𝒂𝒏 𝒄𝒐𝒔 𝑳 𝒙

1 𝐿 2 𝐿 𝒏𝝅
where 𝑎0 = 𝐿 ∫0 𝑓 𝑥 𝑑𝑥 and 𝑎𝑛 = 𝐿 ∫0 𝑓 𝑥 𝒄𝒐𝒔 𝒙𝑑𝑥 and 𝑛 is a positive integer.
𝑳

HALL'S MARRIAGE THEOREM: Let 𝑆 be a family of finite sets, where the family may
contain an infinite number of sets and the individual sets may be repeated multiple
times.

A transversal for 𝑆 is a set T and a bijection 𝑓 from 𝑇 to 𝑆 such that for all 𝑡 in 𝑇, 𝑡 is a
member of 𝑓(𝑡). An alternative term for transversal is system of distinct
representatives or "SDR".

The collection 𝑆 satisfies the marriage condition (MC) if and only if for each
subcollection , we have

In other words, the number of sets in each subcollection 𝑊 is less than or equal to the
number of distinct elements in the union over the subcollection 𝑊. Hall's theorem states
that 𝑆 has a transversal (SDR) if and only if 𝑆 satisfies the marriage condition.

HAMEL BASE: A subject 𝐵 of 𝑆 is called a Hamel base (or simply base) for 𝑆 if and only if
𝐵 is a linealy independent set and 𝐿𝑀 𝐵 = 𝑆.

A linear space 𝑆 is called finite-dimensional if an only if 𝑆 has a finite base B, i.e. B is a

finite set which is a Hamel base. If 𝑆 is not finite dimensional it is called infinite
dimensional. 𝐿𝑀 (𝐵) means linear manifold spanned by B.

HAMILTONIAN CIRCUIT: A Hamiltonian circuit in a graph is a simple circuit that passes

through every vertex of the graph. Thus a circuit 𝒗𝟎 𝒗𝟏 𝒗𝟐 . . . 𝒗𝒏−𝟏 𝒗𝟎 in a graph (𝑽, 𝑬) is a
Hamiltonian circuit if and only if every vertex of the graph occurs exactly once in the list
𝒗𝟎 , 𝒗𝟏 , 𝒗𝟐 , . . . 𝒗𝒏−𝟏 .
HAMILTONIAN PATH: A Hamiltonian path in a graph is a path that passes (exactly once)
through every vertex of the graph. Thus a path 𝒗𝟎 𝒗𝟏 𝒗𝟐 . . . 𝒗𝒏 in a graph (𝑽, 𝑬) is a
Hamiltonian path if and only if 𝑽 = {𝒗𝟎 , 𝒗𝟏 , 𝒗𝟐 , . . . 𝒗𝒏 }. A Hamiltonian path passes can
have no repeated vertices (since it is a path) and therefore passes through each vertex
of the graph exactly once.

HAMILTONIAN SYSTEM: A system of ordinary differential equations in

unknowns and 𝒒 = 𝒒𝟏 , 𝒒𝟐 , … , 𝒒𝒏 , of the form
𝒅𝒑𝒍 𝝏𝑯 𝒅𝒒𝒍 𝝏𝑯
=− , = : 𝒊 = 𝟏, 𝟐, … . , 𝒏
𝒅𝒕 𝝏𝒒𝒊 𝒅𝒕 𝝏𝒑𝒊
where 𝑯 is some function of (𝒑, 𝒒, 𝒕), known as the Hamilton function, or Hamiltonian,
of the system (1). A Hamiltonian system is also said to be a canonical system and in the
autonomous case it may be referred to as a conservative system, since in this case the
function 𝑯 (which often has the meaning of energy) is a first integral (i.e. the energy is
conserved during motion).
HAMILTON-JACOBI THEORY:A branch of classical variational calculus and analytical
mechanics in which the task of finding extremals (or the task of integrating a
Hamiltonian system of equations) is reduced to the integration of a first-order partial
differential equation — the so-called Hamilton–Jacobi equation. The fundamentals of
the Hamilton–Jacobi theory were developed by W. Hamilton in the 1820s for problems
in wave optics and geometrical optics. In 1834 Hamilton extended his ideas to problems
in dynamics, and C.G.J. Jacobi (1837) applied the method to the general problems of
classical variational calculus.
HAMILTON OPERATOR: A symbolic first-order differential operator, used for the
notation of one of the principal differential operations of vector analysis. In a
rectangular Cartesian coordinate system 𝒙 = 𝒙𝟏 , 𝒙𝟐 , … … … , 𝒙𝒏 with unit
vectors 𝒆𝟏 , 𝒆𝟐 , … … … , 𝒆𝒏 , the Hamilton operator has the form
𝒏
𝝏
𝛁= 𝒆𝒋
𝝏𝒙𝒋
𝒋=𝟏

The application of the Hamilton operator to a scalar function 𝑓, which is understood as

multiplication of the "vector" ∆ by the scalar 𝑓(𝑥), yields the gradient of 𝑓:
𝑛
𝜕𝑓
𝑔𝑟𝑎𝑑 𝑓 = ∆𝑓 = 𝒆𝒋
𝜕𝑥𝒋
𝑗 =1

HAMILTON-OSTROGRADSKI PRINCIPLE: A general integral variational principle of

classical mechanics, established by 𝑊. Hamilton for holonomic systems restricted by
ideal stationary constraints, and generalized by M.V. Ostrogradski to non-stationary
geometrical constraints. According to this principle, in a real motion of the system acted
upon by potential forces,
𝑡1 𝑡1

𝑆= (𝑇 − 𝑈)𝑑𝑡 = 𝐿𝑑𝑡
𝑡0 𝑡0

has a stationary value as compared with near, kinetically-possible, motions, with initial
and final positions of the system and times of motion identical with those for the real
motion. Here, 𝑇 is the kinetic energy, 𝑈 is the potential energy and 𝐿 = 𝑇 − 𝑈 is the
Lagrange function of the system.
HAMILTON’S PRINCIPLE OF LEAST ACTION: A particle moves in a conservative field in
𝑡
such a way that ∫𝑡 2 𝑇 − 𝑉 𝑑𝑡 is a minimum, where T is the total Kinetic energy and V is
1

the potential energy of the particle.

HAMMING CODE WEIGHT: Let 𝑥 ∇ 𝐹𝑞𝑛 . The Hamming weight of 𝑥, denoted 𝑤𝑡(𝑥) is
defined to be the number of coordinates that are not zero. That is, 𝑤𝑡(𝑥) = 𝑑(𝑥, 0).
HAMMING DISTANCE: Let 𝑥 = 𝑥1 , . . . , 𝑥𝑛 and 𝑦 = 𝑦1 , . . . , 𝑦𝑛 . Then, for every 𝑖, we
define
1 𝑖𝑓 𝑥𝑖 ≠ 𝑦𝑖
𝑑 𝑥𝑖 , 𝑦𝑖 =
0 𝑖𝑓 𝑥𝑖 = 𝑦𝑖
and we define

𝑑(𝑥, 𝑦) = 𝑑 𝑥𝑖 , 𝑦𝑖

This distance is called Hamming distance. Note that that the Hamming distance is not
dependent on the actual values of 𝑥𝑖 and 𝑦𝑖 but only if they are equal to each other or
not equal.
HANKEL MATRIX: A matrix whose entries along a parallel to the main anti-diagonal are
equal, for each parallel. Equivalently, 𝑯 = (𝒉𝒊,𝒋 ) is a Hankel matrix if and only if there
exists a sequence 𝒔𝟏 , 𝒔𝟐 , … , such that 𝒉𝒊,𝒋 = 𝒔𝒊+𝒋−𝟏 . If 𝒔𝒌 are square matrices, then 𝑯 is
referred to as a block Hankel matrix. Infinite Hankel matrices are associated with the
representation of Hankel operators acting on the Hilbert space of square summable
complex sequences.
HANSON- HUARD STRICT CONVERSE DUALITY THEOREM: Let X 0 be an open set in Rn
and let 𝜃 and 𝑔 be differentiable on X 0 . Let x, u be a solution of the dual
(maximization) problem and let 𝜃 and 𝑔 be convex at 𝑥 if either.

 Ψ(x, u) is twice continuously differentiable at 𝑥 and the 𝑛 × 𝑛 Hessian matrix

∆2x Ψ(x, u) is non singular, or
 There exists an open set ⋀ ⊂ 𝑅 𝑚 containing 𝑢 and an 𝑛-dimensional
differentiable vector function 𝑒(𝑢) on ⋀ such that x = 𝑒(u) and 𝑒 u ∇
X 0 , ∆x Ψ(x, u) 𝑥=𝑒(𝑢) = 0 for 𝑢 ∇ ⋀

then 𝑥 solves the (primal) minimization problem (MP) and θ x = Ψ(x, u).

HARDY INEQUALITY: If 𝑝 > 1, 𝑎𝑛 ≥ 0 and 𝑥𝑛 = 𝑎1 + 𝑎2 + − − +𝑎𝑛 , then

∞ ∞
𝑥𝑛 𝑝 𝑝 𝑝
< 𝑎𝑛𝑝
𝑛 𝑝−1
𝑛=1 𝑛=1

𝑝 𝑝
except when all the 𝑎𝑛 are zero. The constant in this inequality is best possible.
𝑝−1

HARDY–LITTLEWOOD INEQUALITY: It states that if 𝑓 and 𝑔 are nonnegative

measurable real functions vanishing at infinity that are defined on 𝑛-
dimensional Euclidean space 𝑹𝒏 then

where 𝑓 * and 𝑔∗ are the symmetric decreasing rearrangements of 𝑓(𝑥) and 𝑔(𝑥),
respectively.

HARDY-RAMANUJAN THEOREM: For any integer 𝒏 ≥ 𝟐 , let 𝝎(𝒏) denote the number of
distinct prime factors of 𝒏. The Hardy–Ramanujan theorem states that the
function 𝝎(𝒏) has normal order 𝒍𝒐𝒈 𝒍𝒐𝒈 𝒏 in the sense that, given any 𝜺 > 0, almost all
positive integers 𝒏 satisfy
𝝎 𝒏 − 𝐥𝐨𝐠 𝐥𝐨𝐠 𝒏 ≤ 𝜺 𝐥𝐨𝐠 𝐥𝐨𝐠 𝒏

Here, "almost all" means that the number of positive integers 𝑛 ≤ 𝑥 for which the
indicated property holds is asymptotically equal to 𝑥 as 𝑥 → ∞. A stronger, and best-
possible, version of this result shows that, given any real-valued function 𝜑(𝑛) tending
to infinity as 𝑛 → ∞, almost all positive integers 𝑛 satisfy
 𝜔 𝑛 − log log 𝑛 ≤ 𝜑(𝑛) 𝑙𝑜𝑔⁡𝑙𝑜𝑔⁡𝑛

HARDY'S INEQUALITY: Hardy's inequality is an inequality in mathematics, named

after G. H. Hardy. It states that if is a sequence of non-negative real
numbers which is not identically zero, then for every real number p > 1 one has

HARDY'S THEOREM: Let 𝑓 be a holomorphic function on the open ball centered at zero
and radius 𝑅 in the complex plane, and assume that 𝑓 is not a constant function. If one
defines

for 0 < 𝑟 < 𝑅 then this function is strictly increasing and logarithmically convex.

HARMONIC ANALYSIS: A name given to a branch of mathematics and to a mathematical

method. Harmonic analysis as a branch of mathematics is usually understood to include
the theory of trigonometric series; Fourier transforms; almost-periodic
functions; Dirichlet series; approximation theory; abstract harmonic analysis; and
certain other related mathematical disciplines. The method consists in reducing certain
problems from various fields of mathematics to problems in harmonic analysis, which
are then solved.
HARMONIC COORDINATES: Coordinates in which the metric tensor 𝒈𝒊𝒌 satisfies the
condition
𝜕
(|𝑔| − − 𝒈𝒊𝒌 ) = 0,
𝜕𝑥𝑘

where 𝑔 is the determinant defined by the components of the tensor 𝒈𝒊𝒌 In several cases
use of harmonic coordinates leads to a considerable simplification of the calculations:
an example is the derivation of the equations of motion in general relativity.
HARMONIC FUNCTION: A function 𝑢 𝑥, 𝑦 is called harmonic function if first and second
order partial derivatives of 𝑢 are continuous and 𝑢 satisfies Laplace’s equations ∆2 V =
0.
HARMONIC SEQUENCE: A sequence of numbers is a harmonic sequence if the
reciprocals of the terms form an arithmetic sequence.
HCF: It is the abbreviation for highest common factor of two numbers. (Also written as
gcd).
HEAVISIDE UNIT FUNCTIONS (OR THE UNIT STEP FUNCTION): The Heaviside unit step
function by H(x) is defined as

1 if x ≥ 𝑎
H x =
0 if x < 𝑎

If 𝑎 > 0, then

1 if x ≥ 𝑎
H x−𝑎 =
0 if x < 𝑎

HEIGHT OF A FUZZY SET: The height of a fuzzy set is the largest membership value
attained by any point.
HELICOID: A surface generated by a curve which is simultaneously rotated about a fixed
axis and translate in the direction of the axis with a velocity proportional to the velocity
of rotation, is called helicoids.

HELMHOLTZ THEOREM OF CLASSICAL MECHANICS: Let

be the Hamiltonian of a one-dimensional system, where

is the kinetic energy and

is a "𝑈-shaped" potential energy profile which depends on a parameter 𝑉. Let

denote the time average. Let
Then

HEREDITARY PROPERTY: A property of a space is hereditary if each of its subspaces

possesses this property. Being countable is a hereditary property.

HERMITE DIFFERENTIAL EQUATION: The differential equation

d2 y dy
2
− 2𝑥 + 2λy = 0
dx dx

where λ is a constant, is called Hermite’s differential equation.

HERMITE–HADAMARD INEQUALITY: It is also called Hadamard's inequality, states that

if a function ƒ : [𝑎, 𝑏] → 𝑹 is convex, then the following chain of inequalities hold:

th
HERMITIAN MATRIX: A square matrix A = aij is said to be Hermitian if the i, j
th
element A is equal to the conjugate complex of the j, i element A i.e., if aij = aji for all
i and j.

1 2 − 3i 3 + 4i
a b + ic
For example, , 2 + 3i 0 4 − 5i are Hermitian matrices. If A is a
b − ic d
3 − 4i 4 + 5i 2
Hermitian matrix, then

aii = aii (by definition).

∴ aii is real for all i.

Thus every diagonal element of a Hermitian matrix must be real.

A Hermitian matrix over the field of real number is nothing but a real symmetric matrix.
Obviously, a necessary and sufficient condition for a matrix A to be Hermitian is that A
= A°.
HERMITIAN OPERATOR OR SELF-ADJOINT OPERATOR: An operator 𝑇: 𝐻 → 𝐻 is a
Hermitian operator or self-adjoint operator if
𝑇 = 𝑇 ∗ , 𝑖. 𝑒. ⟨ 𝑇𝑥, 𝑦 ⟩ = ⟨ 𝑥, 𝑇𝑦 ⟩ 𝑓𝑜𝑟 𝑎𝑙𝑙 𝑥, 𝑦 ∇ 𝐻.

For any bounded operator 𝑇 operators 𝑇1 = 1/2(𝑇 ± 𝑇 ∗ ), 𝑇 ∗ 𝑇 𝑎𝑛𝑑 𝑇𝑇 ∗ are Hermitians.

HESSIAN MATRIX: The Hess or Hessian matrix of a multivariable function is the

symmetric matrix of second partial derivatives. If 𝑓 (𝑥, 𝑦, 𝑧) is a function of three
variables, its Hessian matrix is
𝑓𝑥𝑥 𝑓𝑥𝑦 𝑓𝑥𝑧
𝑓𝑦𝑥 𝑓𝑦𝑦 𝑓𝑦𝑧
𝑓𝑧𝑥 𝑓𝑧𝑦 𝑓𝑧𝑧
HEWITT- SAVAGE ZERO-ONE LAW: Every symmetric event concerning a sequence of
independent and identically distributed random variables has probability 0 or 1.

HIGHER ORDER DERIVATIVES: If a function 𝑓(𝑧) is analytic within and on a closed

contour 𝐶 and a is any point within 𝐶 then derivatives of all orders are analytic and are
given by

𝑛! 𝑓 𝑧 𝑑𝑧
𝑓 (𝑛) 𝑎 = ∫𝐶
2𝜋𝑖 (𝑧 − 𝑎)𝑛+1

HILBERT, DAVID: David Hilbert (1862 to 1943) was a mathematician whose work
included a modern, rigorous, axiomatic development of geometry.
HILBERT'S BASIS THEOREM: If is a Noetherian ring, then 𝑅[𝑋] is a Noetherian ring.

HILBERT–SCHMIDT NORM: We define Hilbert–Schmidt norm of a Hilbert–Schmidt

2 ∞ 2
operator 𝐴 by 𝐴 𝐻𝑆 = 𝑛=1 𝐴𝑒𝑛 . The set of Hilbert–Schmidt operators with this
norm is a Hilbert space.

HILBERT–SCHMIDT OPERATOR: Let 𝑇: 𝐻 → 𝐾 be a bounded linear map between two

Hilbert spaces. Then 𝑇 is said to be Hilbert–Schmidt operator if there exists an
∞ 2
orthonormal basis in 𝐻 such that the series 𝑘=1 ||𝑇 𝑒𝑘 || is convergent.

Note that

 Let 𝑇: 𝑙2 → 𝑙2 be a diagonal operator defined by 𝑇𝑒𝑛 = 𝑒𝑛 /𝑛, for all 𝑛 ≥ 1.

2
Then 𝑇𝑒𝑛 2
= 𝑛−2 = 𝜋 6 is finite.
  The identity operator 𝐼𝐻 is not a Hilbert–Schmidt operator, unless 𝐻 is finite
dimensional.
 All Hilbert–Schmidt operators are compact.
 Integral operator is Hilbert–Schmidt.

HILBERT SPACE: An (abstract) Hilbert space is a linear inner product space which is
complete with respect to the norm generated by its inner product.

HINGE THEOREM: It states that if two sides of one triangle are congruent to two sides of
another triangle, and the included angle of the first is larger than the included angle of
the second, then the third side of the first triangle is longer than the third side of the
second triangle.
HISTOGRAM: A histogram is a bar diagram where the horizontal axis shows different
categories of values and the height of each bar is related to the number of observations
in the corresponding category. If all categories are the same width, then the height of
each bar is proportional to the number of observations in the category. If the categories
are of unequal width, then the height of the bar is proportional to the number of
observations in the category divided by the width of the category
HÖLDER’S INEQUALITY: For 1 < 𝑝 < ∞, let 𝑞 ∇ (1, ∞) be such that
1 1
+ = 1.
𝑝 𝑞
For 𝑛 ≥ 1 and 𝑢, 𝑣 ∇ 𝐾 𝑛 , we have that
1 1
𝑛 𝑛 𝑝 𝑛 𝑞
𝑝 𝑞
𝑢𝑖 𝑣𝑖 ≤ 𝑢𝑖 𝑣𝑖
𝑖=1 𝑖=1 𝑖=1

HÖLDER’S INEQUALITY FUNCTIONAL ANALYSIS): Let 𝑝 𝑎𝑛𝑑 𝑞 be real numbers greater

than 1, with the added property that

1 1
+ =1
𝑝 𝑞

Then for any 𝑓 = 𝑓 1 , 𝑓 2 , … ∇ 𝑙 𝑃

And = 𝑔 1 , 𝑔 2 , … ∇ 𝑙 𝑞 ,
 ∞ ∞ 1/𝑝 ∞ 1/𝑝
𝑃 𝑞
𝑓 𝑘 𝑔(𝑘) ≤ 𝑓(𝑘) . 𝑔(𝑘)
𝑘=1 𝑘=1 𝑘=1

HOLOMORPHIC FUNCTION: Functions f that satisfy (𝜕/𝜕𝑧)𝑓 ≡ 0 are the main concern
of complex analysis. A continuously differentiable function 𝑓 ∶ 𝑈 → 𝐶 defined on an
open subset 𝑈 of 𝐶 is said to be holomorphic if 𝜕𝑓/ 𝜕𝑧 = 0 at every point of 𝑈. Note
that this last equation is just a reformulation of the Cauchy-Riemann equations.
HOMEOMORPHISM: (a) A continuous bijective map 𝑓 ∶ 𝑋 → 𝑌, such that 𝑓 −1 ∶ 𝑌 → 𝑋
is also continuous is called a homeomorphism (or a bicontinuous bijection) and denoted
by 𝑓 ∶ 𝑋 ≅ 𝑌 . Two spaces 𝑋, 𝑌 are homeomorphic, written 𝑋 ≅ 𝑌 , if there is a
homeomorphism 𝑓 ∶ 𝑋 ≅ 𝑌 .
(b) Let (𝐴,∗) and (𝐵,∗) be semigroups, monoids or groups. A function 𝑓: 𝐴 → 𝐵 from 𝐴
to 𝐵 is said to be a homomorphism if 𝑓(𝑥 ∗ 𝑦) = 𝑓(𝑥) ∗ 𝑓(𝑦) for all elements 𝑥 and 𝑦
of 𝐴. Let 𝑞 be an integer, and let 𝑓: 𝑍 → 𝑍 be a the function from the set of integers to
itself defined by 𝑓(𝑛) = 𝑞 𝑛 for all integers 𝑛. Then 𝑓 is a homomorphism from the
group (𝑍, +) to itself, since 𝑓(𝑚 + 𝑛) = 𝑞(𝑚 + 𝑛) = 𝑞 𝑚 + 𝑞 𝑛 = 𝑓(𝑚) + 𝑓(𝑛) for
all integers 𝑚 and 𝑛. Example Let 𝑅 ∗ denote the set of non-zero real numbers, let 𝑎 be a
nonzero real number, and let 𝑓: 𝑍 → 𝑅 ∗ be the function defined by 𝑓(𝑛) = 𝑎𝑛 for all
integers 𝑚 and 𝑛. Then 𝑓: 𝑍 → 𝑅 ∗ is a homomorphism from the group (𝑍, +) of integers
under addition to the group (𝑅 ∗ ,×) of non-zero real numbers under multiplication,
since 𝑓(𝑚 + 𝑛) = 𝑎 𝑚 +𝑛 = 𝑎𝑚 𝑎𝑛 = 𝑓(𝑚)𝑓(𝑛) for all integers 𝑚 and 𝑛.
HOMOGENEOUS LINEAR EQUATIONS: Suppose

a11 x1 + a12 x2 + ⋯ + a1n xn = 0,

a21 x1 + a22 x2 + ⋯ + a2n xn = 0,
…………………………………… …1
……………………………………
am1 x1 + am2 x2 + ⋯ + amn xn = 0

is a system of m homogeneous equation in 𝑛 unknowns 𝑥1 , 𝑥2 , … , 𝑥𝑛 . Let A=

𝑥1 0
a11 a12 … a1n 𝑥2 0
a21 a22 … a2n 𝑥
X= ⋯ 3 0
O= ⋯ where A, X, O are m × n, n × 1, m × 1
… … … …
am1 am2 … amn m×n, ⋯ ⋯
𝑥𝑛 n×1, 0 m×1,
matrices respectively. Then obviously we can write the system of equations (1) in the
form of a single matrix equation
𝐴𝑋 = 0. … 2

The matrix 𝐴 is called the coefficient matrix of the system of equation (1).

Obviously x1 = 0, x2 = 0, … , xn = 0 i.e., X = 0 is a solution of (1). It is a trivial (self-

obvious) solution of (1).

Again suppose x1 and x2 are two solutions of (2). Then their linear combination,
k1 x1 + k 2 x2 , where k1 and k 2 are any arbitrary numbers, is also a solution of (2).

We have

𝐴 k1 x1 + k 2 x2 = k1 (AX1 ) + k 2 (AX2 ) = k1 0 + k 2 = O[∴ AX1 = O and AX 2 = O]

Hence k1 X1 + k 2 X2 is also a solution of (2).

Therefore the collection of all solution of the system of equation AX = O forms a sub-
space of the n variable, AX = O, is n − r , where r is the rank of the matrix A.

HOMOLOGY GROUP: The 𝑖th homology group 𝐻𝑖 (𝐶 ∗ ) of the complex 𝐶 ∗ is defined to be

the quotient group 𝑍𝑖 (𝐶 ∗ )/𝐵𝑖 (𝐶 ∗ ), where 𝑍𝑖 (𝐶 ∗ ) is the kernel of 𝜕𝑖 ∶ 𝐶𝑖 → 𝐶𝑖−1 and
𝐵𝑖 (𝐶 ∗ ) is the image of 𝜕𝑖+1 ∶ 𝐶𝑖+1 → 𝐶𝑖 . Note that if the modules 𝐶 ∗ occuring in a chain
complex 𝐶 ∗ are modules over some unital ring 𝑅 then the homology groups of the
complex are also modules over this ring 𝑅.
HOMOMORPHISMS OF LIE GROUPS: Let 𝐺, 𝐻 be Lie groups. A map 𝜑: 𝐺 → 𝐻 which is
smooth and a homomorphism in the group theoretical sense, that is, which satisfies
𝜑(𝑥𝑦) = 𝜑(𝑥)𝜑(𝑦) ∀𝑥, 𝑦 ∇ 𝐺,

is called a homomorphism of Lie groups.

HOOKE’s LAW: Theory of elasticity deals with a study of the behavior of those
substances that possess the property of recovering their size and shape, when the forces
producing deformations are removed. The elastic response of a body to an applied load
is analyzed best through the introduction of the stress tensor. Deformations produced
are naturally relate to the applied loads and therefore to the stress tensor. These
relations are the stress- strain relations or the constitutive equations of the material.
In 1676, Robert Hooke, who as a result of his experiments with metallic rods under
axially applied tensile loads, concluded that “the extension is proportional to the force”,
a relation known as Hooke’s Law.

Mathematically it is expresses as T=𝐸ℯ, until the stress reaches the proportional limit.
The proportional limit means the stress at which the linear relationship between stress
and strains ceases. Its value, however, is not easily measured. The constant of
proportionality E is known as Young’s modulus of elasticity.

The stress T is assumed to vanish when ℯ vanishes. Since T and ℯ uniquely related, there
cannot be energy dissipation by the stress component during a loading and unloading
cycle. Such a solid, which allows linear stress-strain law for deformation without any
dissipation of energy, is called a linear elastic solid or Hookean solid.
HOROSPHERE: Horosphere is a level set of Busemann function.
HURWITZ'S THEOREM (COMPLEX ANALYSIS): Let {𝑓𝑘 } be a sequence of holomorphic
functions on a connected open set 𝐺 that converge uniformly on compact subsets of 𝐺 to
a holomorphic function 𝑓. If 𝑓 has a zero of order 𝑚 at 𝑧0 then for every small enough
𝜌 > 0 and for sufficiently large 𝑘 ∇ 𝑵 (depending on 𝜌), 𝑓𝑘 has precisely 𝑚 zeroes in
the disk defined by |𝑧 − 𝑧0 | < 𝜌, including multiplicity. Furthermore, these zeroes
converge to 𝑧0 as 𝑘 → ∞.

HURWITZ'S THEOREM (NUMBER THEORY): For every irrational number 𝜉 there are
infinitely many relatively prime integers 𝑚, 𝑛 such that

The hypothesis that 𝜉 is irrational cannot be omitted. Moreover the constant 5 is the
1+ 5
best possible; if we replace 5 by any number 𝐴 > 5 and we let 𝜉 = (the golden
2

ratio) then there exist only finitely many relatively prime integers 𝑚, 𝑛 such that the
formula above holds.

HYPERBOLA: A hyperbola is the set of all points in a plane such that the difference
between the distances to two fixed points is a constant. A hyperbola has two branches
that are mirror images of each other. Each branch looks like a misshaped parabola. The
general equation for a hyperbola with center at the origin is
 𝑥2 𝑦2
− =1
𝑎2 𝑏 2
The two diagonal lines are called asymptotes, which are determined by the equations
𝑏𝑦 + 𝑎𝑥 = 0 = 𝑏𝑦 − 𝑎𝑥.

HYPERBOLIC POINTS (DIFFERENTIAL GEOMETRY): The points on the surface at which

the Gaussian curvature 𝜅 is negative i.e, 𝐿𝑁 − 𝑀² < 0 are called hyperbolic points. In
this case principal curvatures at the points are of opposite signs.

HYPERGEMETRIC FUNCTION: The symbol α r is defind by

𝑇 α+r
α r = α α + 1 (α + 2)……….. α + r − 1 = (where r is a+ve integer)
𝑇 α

and α 0 = 1.

The general hypergeometric function mFn (α1 , α2 …………..αm ; β1 , β2 ………,βn ; x) is defined

by

∞ α 1 r α 2 r … α m r xT
mFn (α1 , α2 …………..αm ; β1 , β2 ………,βn ; x)= r=0 β
1 r β2 r … βn r r !

Another notation often used for general hypergeometric function is

α1 , α2 … αm ;
mFn = β , β … β 𝑥
1 2 n

HYPERGEOMETRIC SERIES: The series

ab a a + 1 b(b + 1) 2 a a + 1 a + 2 b b + 1 (b + 2) 3
1+ x+ x + x +
1! c 2! c(c + 1) 3! c c + 1 (c + 2)
is called the hypergeometric series.

It can be shown that the series converges absolutely if x < 1 and if x = 1, and series
converges absolutely if c − a − b > 0 or c > 𝑎 + 𝑏. Hypergeometric series is frequently
used in connection with the theory of Spherical Harmonics.

HYPERPLANE: A hyperplane is defined as the set of points satisfying

𝐶1 𝑥1 + 𝐶2 𝑥2 + ⋯ + 𝐶𝑛 𝑥𝑛 =Z (not all 𝐶𝑖 = 0)

or Cx=Z for prescribing values of 𝐶1 , 𝐶2 , … , 𝐶𝑛 and Z.

𝐶
The vector C is called vector normal to hyperplane and ± are called unit normals.
𝐶

A hyper plane divides the whole space 𝐸 𝑛 into three mutually disjoint sets given by

𝑋1 = 𝑥: 𝐶𝑥 > 𝑍
𝑋2 = 𝑥: 𝐶𝑥 = 𝑍

𝑋3 = 𝑥: 𝐶𝑥 < 𝑍

The sets 𝑋1 and 𝑋3 are called open spaces.

The sets 𝑥: 𝐶𝑥 ≤ 𝑍 and 𝑥: 𝐶𝑥 ≥ 𝑍 are called closed half spaces.

HYPER SPHERE: A hyper sphere in 𝐸 𝑛 with the centre at a and radius 𝜀 > 0 is defined to
be the set of points

𝑋 = 𝑥: 𝑥 − 𝑎 = 𝜀

i.e., the equation of hypersphere is 𝐸 𝑛 is

2 2 2
𝑥1 − 𝑎1 + 𝑥2 − 𝑎2 + ⋯ + 𝑥𝑛 − 𝑎𝑛 = 𝜀2

where 𝑎 = 𝑎1 , 𝑎2 , … , 𝑎𝑛 , 𝑋 = 𝑥1 , 𝑥2 , … , 𝑥𝑛

HYPOTHESIS: A hypothesis is a proposition that is being investigated; it has yet to be

proved.
HYPOTHESIS TESTING: A situation often arises in which a researcher needs to test a
hypothesis about the nature of the world. Frequently it is necessary to use a statistical
technique known as hypothesis testing for this purpose. The hypothesis that is being
tested is termed the null hypothesis. The other possible hypothesis, which says “The
null hypothesis is wrong,” is called the alternative hypothesis.

I
IDEA: IDEA (International Data Encryption Algorithm) is an encryption
algorithm developed at ETH in Zurich, Switzerland. It uses a block cipher with a 128-
bit key, and is generally considered to be very secure. It is considered among the best
publicly known algorithms. In the several years that it has been in use, no practical
attacks on it have been published despite of a number of attempts to find some.

IDEAL CLASS: Suppose that 𝐾 is an algebraic number field. Two fractional ideals 𝐴 and 𝐵
in 𝐾 are said to be equivalent, denoted by 𝐴 ∼ 𝐵, if there is a principal fractional ideal
⌌𝑤⌍ ≠ ⌌0⌍ such that 𝐴 = ⌌𝑤⌍𝐵. The equivalence of non-zero fractional ideals in an
algebraic number field is an equivalence relation. The equivalence classes are called the
ideal classes. Furthermore, all non-zero principal fractional ideals are equivalent to each
other. They form the principal class.

IDEAL FLUID OR PERFECT FLUID (FRICTIONLESS, HOMOGENEOUS AND

INCOMPRESSIBLE): The ideal fluid is one which is incapable of sustaining any tangential
stress or action in the form of a shear but the normal force (pressure) acts between the
adjoining layers of fluid. The pressure at every point of an ideal fluid is equal in all
directions, whether the fluid is at rest or in motion. This theory defines some concepts
of the flow such as wave motion. The lift and the induced drag of an airfoil etc., but it
fails to define the phenomena such as skin friction, drag of a body etc.

IDEALS: A subset of a ring 𝐴 is called a left (right) ideal of 𝐴 it is submodule of the left
(right) 𝐴- module of 𝐴. In otherwords, a left (right) ideal 𝐽 of 𝐴 is an additive subgroup
of 𝐴 such that 𝐴𝐽 ⊂ 𝐽(𝐽𝐴 ⊂ 𝐽). Under the operations induced from 𝐴, 𝐽 is a ring
(however, 𝐽 is not necessarily unitary). A subset if 𝐴 is called a two sided ideal or simple
an ideal of 𝐴 if it is a left and right ideal.

For an ideal 𝐽 of a ring 𝐴, we define a relation 𝑅 in 𝐴 by 𝑎𝑅𝑏 ⟺ 𝑎 − 𝑏 ∇ 𝐽. Then 𝑅 is an

equivalence relation that is compatible with the operations of 𝐴. Each equivalence class
is called a residue class modulo 𝐽, and the quotient ring 𝐴/𝑅 is denoted by 𝐴/𝐽 and
called the residue (class) rise (or factor ring) modulo 𝐽. If it is a field, it is called a
residue (class) field. conversely, given an equivalence relation 𝑅 that is compatible with
the operations of 𝐴, the equivalence class of 0 forms an ideal 𝐽 of 𝐴 and the equivalence
relation defined by 𝐽 coincides with 𝑅.

A left (right) ideal of a ring 𝐴 is said to be maximal if it is not equal to 𝐴 and is properly
contained in no left (right) ideal of 𝐴 other than 𝐴. Similarly, a left (right) ideal of 𝐴 is
said to be minimal if it is nonzero left (right) ideal of 𝐴.

If 𝑒 is an idempotent element of a unitary ring 𝐴, then 1 − 𝑒 and 𝑒 are orthogonal

idempotent elements, and 𝐴 = 𝐴𝑒 + 𝐴(1 − 𝑒) is the direct sum of left ideals. This is
called Peirce’s left decomposition. Peirce’s right decomposition is defined similarly.

IDEALS IN AN ALGEBRAIC NUMBER FIELD: Let 𝐾 be an algebraic number field, with ring
of integers 𝑂. A subset 𝐴 of 𝑂 is an ideal in 𝐾 if the following condition is satisfied:
If 𝛼, 𝛽 ∇ 𝐴, then 𝜆𝛼 + µ𝛽 ∇ 𝐴 for every 𝜆, µ ∇ 𝑂.
We denote by {0}, the zero ideal.
IDENTITY ELEMENT: If ⟪ stands for an operation (such as addition, subtraction,
multiplication), then the identity element (called 𝑒) for the operation is the number
such that 𝑒 ⟪ 𝑎 = 𝑎, for all a. For example, zero is the identity element for addition,
because 0 + 𝑎 = 𝑎, for all 𝑎 ∇ 𝑅 . 1 is the identity element for multiplication, because
1 × 𝑎 = 𝑎, for all 𝑎 ∇ 𝑅.
IDENTITY MATRIX: An identity matrix is a square matrix with ones along the diagonal
and zeros everywhere else and is denoted 𝐼𝑛 as where n is the order of the matrix 𝐼𝑛 .
For example:
1 0 0
1 0
𝐼2 = 𝑎𝑛𝑑 𝐼3 = 0 1 0
0 1
0 0 1
IDENTITY THEOREM: If 𝑓, 𝑔 ∶ 𝐶 → 𝐶 are holomorphic near 𝑝, and there is a sequence
𝑝 ≠ 𝑧𝑛 → 𝑝 with 𝑓(𝑧𝑛 ) = 𝑔(𝑧𝑛 ), then 𝑓 = 𝑔 near 𝑝.
IMAGE OF A FILTER: If f is a mapping of a set 𝐸 into a set 𝐸0 :
1. The image by 𝑓 of a filter on 𝐸 is a filter on 𝐸0 ,
2. The inverse image by 𝑓 of a filter 𝐹0 on 𝐸0 is a filter on 𝐸 if every set of 𝐹0 meets 𝑓(𝐸).
In particular, if 𝑓 is a mapping of 𝐸 onto 𝐸0 the direct and inverse images of filters are
again filters.
IMAGE OF A POINT: The image of a point is the point that results after the original point
has been subjected to a transformation.
IMMERSED SUBMANIFOLD: Let 𝑀 be an abstract smooth manifold. An immersed
submanifold is a subset 𝑁 ⊂ 𝑀 equipped as a topological space (but not necessarily
with the topology induced from 𝑀) and a smooth structure such that the inclusion map
𝑖: 𝑁 → 𝑀 is smooth with a differential 𝑑𝑖𝑝 : 𝑇𝑝 𝑁 → 𝑇𝑝 𝑀 which is injective for each
𝑝 ∇ 𝑁.
IMMERSION: Let 𝑀 and 𝑁 be abstract manifolds, and let 𝑓: 𝑁 → 𝑀 be a map. Then 𝑓 is
called an embedding/immersion if its image 𝑓(𝑁) can be given the structure of an
embedded/immersed submanifold of 𝑀 onto which 𝑓 is a diffeomorphism.
IMPLICIT DIFFERENTIATION: Implicit differentiation provides a method for finding
derivatives if the relationship between two variables is not expressed as an explicit
function. For example, consider the equation 𝑥 2 + 𝑥𝑦 + 𝑦 2 = 𝑎. This equation defines a
relationship between 𝑥 and 𝑦, but it does not express that relationship as an explicit
function. To find the derivative 𝑑𝑦/𝑑𝑥, take the derivative of 𝑥 2 + 𝑥𝑦 + 𝑦 2 = 𝑎
𝑑 2 𝑑
𝑥 + 𝑥𝑦 + 𝑦 2 = 𝑎
𝑑𝑥 𝑑𝑥
𝑑𝑦 𝑑𝑦
⇒ 2𝑥 + 𝑥 + 𝑦 + 2𝑦 =0
𝑑𝑥 𝑑𝑥
𝑑𝑦 2𝑥 + 𝑦
⇒ =−
𝑑𝑥 𝑥 + 2𝑦
IMPOSSIBILITY OF CIRCLE SQUARING: We cannot construct a square of area equal to a
given circle by a ruler and compass construction.
IMPUTATION: Imputation, in statistics, is the insertion of a value to stand in for missing
data. Analytics programs and methods don't function properly with missing data.
Statistical packages, for example, commonly delete any case with data missing.
INCENTER: The incenter of a triangle is the center of the circle inscribed inside the
triangle. It is the intersection of the three angle bisectors of the triangle.
INCIDENT EDGE: If 𝑣 is a vertex of some graph, if 𝑒 is an edge of the graph, and if
𝑒 = 𝑣 𝑣 ′ for some vertex 𝑣 ′ of the graph, then the vertex 𝑣 is said to be incident to the
edge 𝑒, and the edge e is said to be incident to the vertex 𝑣.
INCIRCLE: The incircle of a triangle is the circle that can be inscribed within the triangle.
INCOMPLETENESS THEOREM: The Incompleteness Theorem is a pair of logical proofs
that revolutionized mathematics. The first result was published by Kurt Gödel (1906-
1978) in 1931 when he was 24 years old. The First Incompleteness Theorem states that
any contradiction-free rendition of number theory contains propositions that cannot be
proven either true or false on the basis of its own postulates. The Second
Incompleteness Theorem states that if a theory of numbers is contradiction-free, then
this fact cannot be proven with common reasoning methods.

INCONSISTENT EQUATIONS: Two equations are inconsistent if they contradict each

other and therefore cannot be solved simultaneously.
INCREASING FUNCTION: A function 𝑓(𝑥) is an monotonically increasing function if
𝑓 (𝑎) ≤ 𝑓 (𝑏) when 𝑎 ≤ 𝑏. A function 𝑓(𝑥) is an strictly increasing function if
𝑓 𝑎 < 𝑓 (𝑏) when 𝑎 < 𝑏 .
INDECOMPOSABLE MODULE: A non-trivial module over a ring 𝑅 is called
indecomposable if whenever 𝑀 = 𝑀1 ⊕ 𝑀2 with 𝑀1 , 𝑀2 submodules then either
𝑀1 = {0} or 𝑀2 = {0}.
INDEPENDENT EVENTS: Two events are independent if they do not affect each other.
For example, the probability that a new baby will be a boy is not affected by the fact that
a previous baby was a girl. Therefore, these two events are independent. If 𝐴 and 𝐵 are
two independent events, the conditional probability that 𝐴 will occur, given that 𝐵 has
occurred, is just the same as the unconditional probability that 𝐴 will occur:
𝑃 𝐴 ∕ 𝐵 = 𝑃(𝐴)
INDEPENDENT VARIABLE: The independent variable is the input number to a function.
In the equation 𝑦 = 𝑓 (𝑥), 𝑥 is the independent variable and 𝑦 is the dependent
variable. An independent variable is a variable that is manipulated to determine the
value of a dependent variable s. The dependent variable is what is being measured in an
experiment or evaluated in a mathematical equation and the independent variables are
the inputs to that measurement. In a simple mathematical equation, for example:
𝑎 = 𝑏/𝑐 the independent variables, 𝑏 and 𝑐 , determine the value of a . Here's a simple
example: A teacher wishes to compare the number of tardy students wearing black with
the number of tardy students wearing pink. In this scenario, clothing color is the
independent variable and the difference in the number of students, categorized by
clothing color, is the dependent variable.
INDISCRETE TOPOLOGY: Let 𝑋 be any non-empty set and ℐ = 𝑋; 𝜑 . Then ℐ is called
the indiscrete topology and (𝑋; ℐ) is said to be an indiscrete space.
INDUCTIVE REASONING: Inductive reasoning is a logical process in which multiple
premises, all believed true or found true most of the time, are combined to obtain a
specific conclusion. Inductive reasoning is often used in applications that involve
prediction, forecasting, or behavior.

Inductive reasoning is, unlike deductive reasoning, not logically rigorous. Imperfection
can exist and inaccurate conclusions can occur, however rare; in deductive reasoning
the conclusions are mathematically certain. Inductive reasoning is sometimes confused
with mathematical induction, an entirely different process. Mathematical induction is a
form of deductive reasoning, in which logical certainties are "daisy chained" to derive a
general conclusion about an infinite number of objects or situations.

INFINITE FOURIER COSINE TRANSFORM: Let 𝑓 𝑡 be a function defined and piecewise

∞
continuous on 0, ∞ and is absolutely convergent on 0, ∞ i.e. ∫𝟎 𝑓 𝑡 𝑑𝑡 𝑐𝑜𝑛𝑣𝑒𝑟𝑔𝑒𝑠.

Then Fourier cosine transform of 𝑓 𝑡 ,denoted by 𝐹𝑐 𝑓 𝑡 and is defined as

∞
𝐹𝑐 𝑓 𝑡 = 𝑓 𝑡 cos 𝑠 𝑡 𝑑𝑡 = 𝐺𝑐 𝒔
𝟎

and the function 𝒇 𝒕 is inverse fourier cosine transform of 𝐺𝑐 𝑠 𝑜𝑟 𝐹𝑐 𝑓 𝑡 and

defined as

∞
2
𝑓 𝑡 = 𝐺𝑐 𝒔 cos s t ds
𝜋 𝟎

(2) is also known as Inversion Formula for Fourier sine transform.

INFINTE FOURIER SINE TRANSFORM: Let 𝑓 𝑡 be a function defined and piecewise

∞
continuous on 0, ∞ and is absolutely convergent on 0, ∞ i.e. ∫𝟎 𝑓 𝑡 𝑑𝑡 converges

Then Fourier sine transform of 𝑓 𝑡 , denoted by 𝐹𝑠 𝑓 𝑡 and is defined as

∞
𝐹𝑠 𝑓 𝑡 = 𝑓 𝑡 sin 𝑠 𝑡 𝑑𝑡 = 𝐺𝑠 𝒔
𝟎
and the function 𝒇 𝒕 is inverse fourier sine transform of 𝐺𝑠 𝑠 or 𝐹𝑠 𝑓 𝑡 and defined
as

∞
2
𝑓 𝑡 = 𝐺𝑠 𝒔 sin s t ds
𝜋 𝟎

The second is also known as Inversion Formula for Fourier sine transform.

INFINITE PRODUCT: Let 𝑎𝑛 be a given sequence with terms 𝑎𝑛 ≠ 0(𝑛 = 1,2, … 1). The
∞
dormal infinite product 𝑎1 . 𝑎2 . 𝑎3 … is denoted by n=1 𝑎𝑛 . We call 𝑝𝑛 = 𝑎1 . 𝑎2 … … 𝑎𝑛 its
𝑛𝑡𝑕 partial product. If the sequence 𝑝𝑛 is convergent to a nonzero limit 𝑝, then this
infinite product is said to converge to 𝑝, and 𝑝 is called the value of the infinite product.
We write 𝑎𝑛 = p. if 𝑝𝑛 is not convergent or is convergent to 0, then the infinite
product is called divergent.

INFINITESIMAL: An infinitesimal is a variable quantity that approaches very close to

zero. In calculus ∆𝑥 is usually used to represent an infinitesimal change in 𝑥.
INFINITESIMAL AFFINE DEFORMATION: A transformation of the type

𝑥𝑖′ = 𝛼𝑖0 + 𝛿𝑖𝑗 + 𝛼𝑖𝑗 𝑥𝑗 , 𝑖, 𝑗 = 1,2,3

In which the co-efficient are so small that their products can be neglected in comparison
with the linear terms is called an infinitesimal affine transformations.

INFLECTION POINT: An inflection point on a curve is a point such that the curve is
oriented concave-upward on one side of the point and concave-downward on the other
side of the point. If the curve represents the function 𝑦 = 𝑓(𝑥), then the second
derivative is equal to zero at the inflection point.
INITIAL-BOUNDARY-VALUE PROBLEMS: This time the infinite 𝒙 plane is replace by the
semi-infinite plane, with a boundary condition imposed on 𝒙 = 𝟎.
𝝏𝒖 𝝏𝒖
𝒄 + = 𝟎, 𝒙 > 0, 𝑡 > 0
𝝏𝒙 𝝏𝒕
𝒖 𝒙, 𝟎 = 𝒇(𝒙)
𝒖 𝟎, 𝒕 = 𝒈(𝒕)
.

Consider the case c>0, so that the waves are travelling from small 𝑥 to large 𝑥. The
𝑑𝑥 𝑑𝑡
characteristics are again given by = 𝑐, = 1 ⇒ 𝑥 = 𝑐𝑡 + 𝑥0 . The characteristic
𝑑𝑟 𝑑𝑟

through the origin, splits the 𝑥 − 𝑡 plane into two regions.

𝝏𝒖 𝝏𝒖
Example: 2 𝝏𝒙 + = 𝟎, −∞ < 𝑥 < ∞, 𝑡 > 0 with
𝝏𝒕

1 − 𝑥, 𝑖𝑓 𝑥 ≤ 1
𝑢 𝑥, 0 = , 𝑢 0, 𝑡 = 𝑒 −𝑡 , 𝑡 > 0.
0, 𝑖𝑓 𝑥 > 1

The characteristics are given by 𝑥0 = 𝑥 − 2𝑡

Therefore the plane is split into three regions and the solution is

0, 𝑖𝑓 𝑥 − 2𝑡 > 1,
𝑢 𝑥, 𝑡 = 1 − 𝑥 − 2𝑡 , 𝑖𝑓 0 < 𝑥 − 2𝑡 ≤ 1,
𝑥
𝑒− 𝑡− ,
2 , 𝑖𝑓 𝑥 − 2𝑡 < 0

𝝏𝒖 𝝏𝒖
Example: 2𝑥𝑡 𝝏𝒙 + = 𝑢, −∞ < 𝑥 < ∞, 𝑡 > 0 with the initial condition 𝑢 𝑥, 0 = 𝑥.
𝝏𝒕

The characteristic equations are

𝑑𝑥 𝑑𝑥 𝑑𝑡
= 2𝑥𝑡 ⇒ = 2𝑡𝑑𝑡 ⇒ log 𝑥 = 𝑡 2 + 𝑐 and 𝑑𝑟 = 1 ⇒ 𝑡 = 𝑟.
𝑑𝑟 𝑥

2 2 2
and rearranging we obtain 𝑥 = 𝐴𝑒 𝑡 ⇒ 𝑥 = 𝑥0 𝑒 𝑡 . Thus, 𝑥0 = 𝑥𝑒 −𝑡

𝑑𝑢 𝑑𝑥 𝜕𝑢 𝑑𝑡 𝜕𝑢 𝑑𝑢
On using = 𝑑𝑟 + 𝑑𝑟 , we can write = 𝑢 ⇒ 𝑢 = 𝐴𝑒 𝑟 , where the constant 𝐴 is
𝑑𝑟 𝜕𝑥 𝜕𝑡 𝑑𝑟

actually constant along the characteristic defined by the value of 𝑥0 . Thus, 𝑥0 is really a
function of 𝑥0 and we write 𝑢 = 𝐹(𝑥0 )𝑒 𝑡 , since 𝑡 = 𝑟 and the function 𝐹 is determined
by the initial conditions. At 𝑡 = 0 , we have 𝑥 = 𝑥0 and 𝑢 𝑥0 , 0 = 𝑥0 so that 𝐹 𝑥0 = 𝑥0 .
Replacing 𝑥0 by the above expression involving 𝑥 and 𝑡, we obtain the final solution

2 2
𝑢 𝑥, 𝑡 = 𝑥0 𝑒 𝑡 = 𝑥𝑒 −𝑡 𝑒 𝑡 = 𝑥𝑒 𝑡−𝑡

𝝏𝒖 𝝏𝒖
Example: Consider the initial-boundary-value problem 𝑢2 𝝏𝒙 + = 0, 𝒙 > 0, 𝒕 > 0 with
𝝏𝒕

𝑢 𝑥, 0 = 𝑥, 𝑥 > 0, 𝑢 0, 𝑡 = 0, 𝑡 > 0.

Characteristic equations are

𝑑𝑥 𝑑𝑡 𝑑𝑢
= 𝑢2 and 𝑑𝑟 = 1 ⇒ 𝑡 = 𝑟 and 𝑑𝑟 = 0 ⇒ 𝑢 =constant on characteristic = 𝐹 𝑥0 .
𝑑𝑟

Hence, 𝑥 = 𝑢2 𝑟 + 𝑥0 = 𝑢2 𝑡 + 𝑥0 since 𝑢 is constant along the characteristic and 𝑑𝑥/𝑑𝑟 is

the derivative of 𝑥 along the characteristic curve.

Therefore, we have 𝑥0 = 𝑥 − 𝑢2 𝑡.

Thus, substituting for 𝑥0 into 𝐹(𝑥0 ) we obtain the implicit solution 𝑢 𝑥, 𝑡 = 𝐹(𝑥 − 𝑢2 𝑡),
where 𝐹 is determined by the initial condition, namely, 𝑡 = 0, 𝑥 = 𝑥0 , and 𝑢 = 𝑥0 .

Thus, 𝐹 𝑥0 = 𝑥0 and so 𝑢 = 𝑥 − 𝑢2 𝑡, 𝑥 − 𝑢2 𝑡 > 0.

Squaring both sides, we can manipulate this solution into an explicit solution for 𝑢 in
terms of 𝑥 and 𝑡.

𝑥 𝑥
𝑢2 = 𝑥 − 𝑢2 𝑡, 𝑢2 1 + 𝑡 = 𝑥, 𝑢2 = 1+𝑡 and so the final solution is 𝑢 = , 𝑥 > 0, 𝑡 >
1+𝑡

0.

Some examples of important linear partial differential equations are:

𝜕2𝑢 𝜕2𝑢
= 𝑐 2 𝜕𝑥 2 (one dimensional wave equation),
𝜕𝑡 2
𝜕𝑢 𝜕2𝑢
= 𝑘 2 𝜕𝑦 2 one dimensional heat or diffusion equation ,
𝜕𝑡

𝜕2𝑢 𝜕2𝑢
+ 𝜕𝑦 2 = 0 two dimensional Laplace equation .
𝜕𝑥 2

They are all examples of homogeneous, linear partial differential equations.

An example of an inhomogeneous equation is the two dimensional Poisson equation,
namely

𝜕2𝑢 𝜕2𝑢
+ 𝜕𝑦 2 = 𝑓 𝑥, 𝑦 .
𝜕𝑥 2

If 𝑢1 (𝑥, 𝑦) and 𝑢2 (𝑥, 𝑦) are both solutions to a partial differential equation, then

𝑢 = 𝑐1 𝑢1 𝑥, 𝑦 + 𝑐2 𝑢2 (𝑥, 𝑦) is also a solution if 𝑐1 and 𝑐2 are constants. However, unlike

second order ordinary differential equations where there are two linearly independent
solutions and two arbitrary constants, linear partial differential equations may well
have an infinite number of linearly independent solutions and we may have to add
together solutions involving an infinite number of constants.

INFORMATION THEORY: Information theory is a branch of mathematics that overlaps

into communications engineering, biology, medical science, sociology, and psychology.
The theory is devoted to the discovery and exploration of mathematical laws that
govern the behavior of data as it is transferred, stored, or retrieved.

INITIAL-VALUE PROBLEMS: Here the problem is defined over an infinite range in 𝒙. The
general statement of the problem is

𝝏𝒖 𝝏𝒖
𝒄 𝝏𝒙 + = 𝟎, −∞ < 𝑥 < ∞, 𝑡 > 0. (1)
𝝏𝒕

𝒖 𝒙, 𝟎 = 𝒇 𝒙 . (2)
The initial value problem (IVP) defined in Equations (1) and (2), is also referred to as a
Cauchy problem and the solution is determined uniquely by the single condition at
𝑡 = 0. Also, we have

𝝏𝒖
𝑢 𝑥, 𝑡 = 𝑓 𝑥 − 𝑐𝑡 = 𝑓(𝑥0 ). If 𝑓 is continuously differentiable, then 𝝏𝒙 = 𝑓 ′′ (𝑥 − 𝑐𝑡) and
𝝏𝒖
= −𝑐𝑓 ′ (𝑥 − 𝑐𝑡) are continuous. Thus, 𝑢 = 𝐹 𝑥0 = 𝐹 𝑥 − 𝑐𝑡 is a classical solution of
𝝏𝒕

characteristic curves. If 𝑓 is only piecewise continuous (derivatives are not continuous),

e.g.

1 − (𝑥 − 𝑐𝑡)2 , 𝑖𝑓 0 ≤ (𝑥 − 𝑐𝑡)2 ≤ 1
𝑢= , then 𝑢 𝑥, 𝑡 = 𝑓 𝑥 − 𝑐𝑡 is called a weak
0, 𝑖𝑓 (𝑥 − 𝑐𝑡)2 > 1
solution.
 𝝏𝒖 𝝏𝒖 1
Example: 2 𝝏𝒙 + = 𝟎, −∞ < 𝑥 < ∞, 𝑡 > 0 with the initial condition 𝑢 𝑥, 0 = 1+𝑥 2 .
𝝏𝒕

The characteristic curves are given by the solutions to

𝑑𝑥 𝑑𝑡
= 2 ⇒ 𝑥 = 2𝑟 + 𝑥0 and 𝑑𝑟 = 1 ⇒ 𝑡 = 𝑟. Hence, 𝑥 = 2𝑡 + 𝑥0 ⇒ 𝑥0 = 𝑥 − 2𝑡.
𝑑𝑟

𝑑𝑢
The partial differential equation becomes 𝑑𝑟 = 0 ⇒ 𝑢 = 𝑢(𝑥0 ).

1 1
At 𝑡 = 0, 𝑥 = 𝑥0 and 𝑢 = 1+𝑥 2 . Therefore, 𝑢 𝑥, 𝑡 = 1+(𝑥−2𝑡)2 .
0

𝝏𝒖 𝝏𝒖
Example: − 𝝏𝒙 + = 𝟎, −∞ < 𝑥 < ∞, 𝑡 > 0 with the initial condition
𝝏𝒕

1 − 𝑥 , 𝑖𝑓 𝑥 ≤ 1
𝑢 𝑥, 0 = .
0, 𝑖𝑓 𝑥 > 1

The characteristic curves are given by the solutions to

𝑑𝑥 𝑑𝑡
= −1 ⇒ 𝑥 = −𝑟 + 𝑥0 and 𝑑𝑟 = 1 ⇒ 𝑡 = 𝑟.
𝑑𝑟

Thus, eliminating 𝑟 and re-writing 𝑥0 in terms of 𝑥 and 𝑡, we obtain 𝑥0 = 𝑥 + 𝑡.

𝑑𝑢
The partial differential equation reduces to the ODE 𝑑𝑟 = 0 ⇒ 𝑢 = 𝑢 𝑥0 .

Using the initial condition, 𝑡 = 0, 𝑥 = 𝑥0 , , we have

1 − 𝑥0 , 𝑖𝑓 𝑥0 ≤ 1 1 − 𝑥 + 𝑡 , 𝑖𝑓 𝑥 + 𝑡 ≤ 1
𝑢(𝑥0 , 0) = and so 𝑢 𝑥, 𝑡 = .
0, 𝑖𝑓 𝑥0 > 1 0, 𝑖𝑓 𝑥 + 𝑡 > 1

INJECTIVE FUNCTION: We say that a function 𝑓 ∶ 𝑋 → 𝑌 is injective or one-to-one if

𝑓(𝑥) = 𝑓(𝑦) implies 𝑥 = 𝑦, or equivalently, whenever 𝑥 ≠ 𝑦, then 𝑓(𝑥) ≠ 𝑓(𝑦).
INNER PRODUCT: An inner product (or scalar product) is an assignment of a scalar
(𝑓, 𝑔) to every two vectors 𝑓 𝑎𝑛𝑑 𝑔 of a linear space 𝑆, in such a way that

𝑓, 𝑔 = 𝑔, 𝑓 ,

(𝑓1 + 𝑓2 , g) = (𝑓1 , 𝑔) + 𝑓2 , 𝑔 ,

𝑎𝑓, 𝑔 = 𝑎 𝑓, 𝑔 ,

𝑓, 𝑓 > 0, 𝑖𝑓 𝑓 ≠ 0
INNER PRODUCT SPACE: The linear space 𝑆 in which an inner product has been defined
is called an inner product space.

INTEGRAL CLOSURE OF A RING: Let 𝑇 be a unital commutative ring, let 𝑅 be a unital

subring of 𝑇. The integral closure 𝑅 of 𝑅 in 𝑇 is the subring of 𝑅 consisting of all
elements of 𝑇 that are integral over 𝑅.

INTEGRAL DOMAIN: A unital commutative ring 𝑅 is said to be an integral domain if the

product of any two non-zero elements of 𝑅 is itself non-zero. A field is an integral
domain.

INTEGRAL ELEMENT OVER R: Let 𝑇 be a unital commutative ring, and let 𝑅 be a unital
subring of 𝑇. An element 𝛼 of 𝑇 is said to be integral over 𝑅 if 𝛼 is the root of some
monic polynomial with coefficients in 𝑅.

INTEGRAL HÖLDER INEQUALITY: Let 1 < 𝑝 < ∞, let 𝑞 ∇ (1, ∞) be such that
1/𝑝 + 1/𝑞 = 1. For 𝑓 ∇ 𝐿𝑝 (µ) and 𝑔 ∇ 𝐿𝑞 (µ), we have that 𝑓𝑔 is integrable, and
∫𝑋 𝑓𝑔 𝑑𝜇 ≤ 𝑓 𝑝 𝑔 𝑞

Let 𝑓, 𝑔 ∇ 𝐿𝑝 (µ) and let 𝑎 ∇ 𝐾. Then:

1. ||𝑎𝑓||𝑝 = | 𝑎 | ||𝑓||𝑝 ;

2. || 𝑓 + 𝑔 ||𝑝 ≤ ||𝑓||𝑝 + ||𝑔||𝑝 .

In particular, 𝐿𝑝 is a vector space.

For 1 ≤ 𝑝 < ∞, the collection of equivalence classes 𝐿𝑝 (µ) / ∼ is a vector space, and
|| · ||𝑝 is a well-defined norm on 𝐿𝑝 (µ) / ∼.

INTEGRALLY CLOSED SUBRING: Let 𝑇 be a unital commutative ring, and let 𝑅 be a unital
subring of 𝑇. The subring 𝑅 of 𝑇 is said to be integrally closed in 𝑇 if every element of 𝑇
that is integal over 𝑅 is an element of 𝑅. Let 𝑇 be a unital commutative ring, let 𝑅 be a
unital subring of 𝑇, and let 𝑅 be the integral closure of 𝑅 in 𝑇. R is an integrally-closed
subring of 𝑇. Moreover 𝑅 is integrally closed in 𝑇 if and only if 𝑅 = 𝑅 . An integral
domain is said to be integrally closed if it is integrally closed in its field of fractions.

INTEGRAL OF A SIMPLE FUNCTION: For any summable function 𝑓, we define the

integral of a simple function 𝑓: 𝑋 → ℝ over a measurable set 𝐴 by setting ∫ 𝑓 𝑑𝜇 =

𝐴
 ∞
𝑘=1 𝑡𝑘 𝜇 𝐴𝑘 ∩ 𝐴 . The value of integral of a summable function is independent from
its representation by the sum of indicators over pair-wise disjoint sets.

Let 𝑓, 𝑔: 𝑋 → ℝ be simple summable function, let 𝑎, 𝑏 ∇ ℝ and 𝐴 is a measurable. Then:

1. ∫(𝑎𝑓 + 𝑏𝑔) 𝑑𝜇 = 𝑎 ∫ 𝑓 𝑑𝜇 + 𝑏∫ 𝑓 𝑑𝜇 , that is 𝑆(𝑋) is a linear space;

𝐴 𝐴 𝐴

2. The correspondence 𝑓 → ∫ 𝑓 𝑑𝜇is a linear functional on 𝑆(𝑋);

𝐴

3. The correspondence 𝐴 → ∫ 𝑓 𝑑𝜇 is a charge;

𝐴

4. The function 𝑑1 𝑓, 𝑔 = ∫ 𝑓 𝑥 − 𝑔 𝑥 𝑑𝜇(𝑥) has all properties of the distance

on 𝑆(𝑋) probably except separation.

5. For all 𝐴 ⊂ 𝑋: ∫ 𝑓 𝑥 𝑑𝜇 𝑥 − ∫ 𝑔 𝑥 𝑑𝜇 𝑥 ≤ 𝑑1 (𝑓, 𝑔)

6. If 𝑓 ≤ 𝑔 then ∫ 𝑓 𝑑µ ≤ ∫ 𝑔 𝑑µ, that is integral is monotonic;

7. For 𝑓 ≥ 0 we have ∫ 𝑓 𝑑µ = 0 iff µ( { 𝑥 ∇ 𝑋 ∶ 𝑓(𝑥) ≠ 0 } ) = 0.

INTEGRAND: The integrand is a function that is to be integrated. In the expression

2𝑥+3 2𝑥+3
∫ 𝑑𝑥, the function is the integrand.
𝑥+21 𝑥+21

INTERCEPT: The 𝑥- intercept is the value of 𝑥 where the curve crosses the 𝑥-axis and 𝑦-
intercept of a curve is the value of 𝑦 where it crosses the 𝑦-axis. For the line
𝑐
𝑦 = 𝑚𝑥 + 𝑐, the 𝑥- intercept is − 𝑚 and the 𝑦-intercept is 𝑐.

INTERIOR POINT OF A SET: Let 𝑆 be a subset of 𝑅. A point 𝑝 ∇ 𝑆 is said to be an interior

point of 𝑆 if there exists a nbd (𝑝−∇, 𝑝+∇) of 𝑝 which is entirely contained in 𝑆.

INTERMEDIATE VALUE THEOREM: The intermediate value theorem states the

following: Consider an interval 𝐼 = [𝑎, 𝑏] in the real numbers ℝ and a continuous
function 𝑓 : 𝐼 → ℝ. Then,

𝑓 𝑎 < 𝑢 < 𝑓 𝑏 𝑜𝑟 𝑓 𝑎 > 𝑢 > 𝑓 𝑏 ⇒ ∃ 𝑐 ∇ 𝑎, 𝑏 𝑠. 𝑡. 𝑓(𝑐) = 𝑢.

INTERPOLATION: Interpolation provides a means of estimating the value of a function
for a particular number if we know the value of the function for two other numbers
above and below that number. The general formula of interpolation is
𝑐−𝑎
𝑓 𝑐 =𝑓 𝑎 + 𝑓 𝑏 − 𝑓(𝑎)
𝑏−𝑎
For example, if 𝑓(26) = 655 and 𝑓(29) = 733, then
28 − 26
𝑓 28 = 𝑓 26 + 𝑓 29 − 𝑓(26)
29 − 26
2 2
= 655 + (733 − 655) = 655 + 78 = 655 + 52 = 707
3 3
INTERSECTION OF FUZZY SETS: The intersection of fuzzy sets 𝐴 and 𝐵 is a fuzzy set
𝐴 ∩ 𝐵 in 𝑈 with membership function
𝜇𝐴 ∩ 𝐵 (𝑥) = 𝑚𝑖𝑛[𝜇𝐴 (𝑥), 𝜇𝐵 (𝑥)]
INTERSECTION OF SETS: The intersection of two sets 𝐴 and 𝐵 is the set of all elements
contained in both sets 𝐴 and 𝐵 and is denoted as 𝐴 ∩ 𝐵. For example, the intersection of
the sets 𝐴 = {1,2,3,4,5,6,7,8,9} and 𝐵 = {3,6,9,12,15,18,21} is the set 𝐴 ∩ 𝐵 = {3,6,9}.
INTERSECTION OF SUBSPACES: If 𝑆 and 𝑇 be two subspaces of 𝑉𝑛 , then the vectors
common to both 𝑆 and 𝑇 also constitute a subspace. This subspace is called the
intersection of the subspaces 𝑆 and 𝑇.

INTERVAL: The interval [𝑎, 𝑏] is the set of points between 𝑎 and 𝑏 including both
endpoints 𝑎 and 𝑏 themselves and is called a closed interval. The interval (𝑎, 𝑏) is the set
of points between 𝑎 and 𝑏 excluding both endpoints 𝑎 and 𝑏 themselves and is called an
open interval.
ij ….k
INTRISIC DERIVATIVE (TENSORS): Consider a tensor 𝐴 ab ….p
defined along a curve

𝑥 𝑖 = 𝑥 𝑖 (t) of parameter 𝑡. The intrinsic derivative is defined by the equation

ij ….k
δA ab ….p dx P
= A abij….k
….,p
. It means that the intrinsic derivative of a tensor is a tensor of the
δt dt

same rank and similar character as the original tensor.

δ∅ 𝜕∅ dx i dx i 𝜕∅ d∅
Similarly intrinsic derivative of ∅ is = = ∅, i . Also = . Thus intrinsic
δt 𝜕𝑥 𝑖 dt dt 𝜕𝑡 dt

derivative of ∅ is its total derivative.

INVARIANT: An invariant quantity doesn’t change under specified conditions. For

example, the distance between two points in Euclidian space is invariant if we rotate or
translate the coordinate system used to express those points.
INVENTORY: The inventory can be defined as a stock of goods which is kept for the
future purpose.

TYPES OF INVENTORY:
(i) RAW MATERIALS: The material used in manufacture of the products such as
fuels etc. is called raw material
(ii) PARTLY FINISHED ITEMS: The material which held between manufacturing stage
is called partly finished items.
(iii) FINISHED GOODS: The products which are ready for sale or distribution called
finished goods.
(iv) SPARE PARTS: The spare parts used in the production process but do not become
part of the product.

INVENTORY CONTROL: Inventory control is the process of deciding what and how much
of various terms are to be kept in stock. It also use to determines the time and quantity
of various items to be procured. It is useful to reduce investment in inventories and
ensuring that production process does not suffer at the same time.

INVERSE: If ⟪ represents an operation (such as addition, multiplication, subtraction

etc), and 𝑒 represents the identity element of that operation, then the inverse of a
number 𝑎 is denoted as 𝑎−1 and is the number that satisfies 𝑎 ⟪ 𝑎−1 = 𝑒. For example,
the additive inverse of a number 𝑥 is −𝑥 (also called the negative of x) because
𝑥 + −𝑥 = 0. The multiplicative inverse of 𝑥 is 1/𝑥 (also called the reciprocal of 𝑥)
1
because 𝑥 × 𝑥 = 1(assuming 𝑥 ≠ 0).

INVERSE FINITE FOURIER COSINE TRANSFORM: The above function 𝑓 𝑡 is known as

inverse finite Fourier cosine transform of 𝐹𝑐 𝑓 𝑡 𝑜𝑟𝐹𝑐 𝑠 and is defined as

∞
1 2 𝑐𝑜𝑠 𝜋 𝑠 𝑡
𝐹𝐶−1 𝐹𝐶 𝑓 𝑡 = 𝑓 𝑡 = 𝐹𝑐 0 + 𝐹𝑐 𝑠
ℓ ℓ ℓ
𝑠=1

This formula is obtained from Fourier cosine series

∞
𝑎0 𝑐𝑜𝑠 𝜋 𝑛 𝑡
𝑓 𝑡 = + 𝑎𝑛 .
2 ℓ
𝑛=1

It is also known as inversion formula for finite Fourier cosine transform.

INVERSE FINITE FOURIER SINE TRANSFORM: Let f(t) be a function defined on 0, ℓ and
satisfying Dirichlet’s conditions on 0, ℓ . The finite Fourier sine transform of
𝑓 𝑡 , 𝑜 < 𝑡 < 𝑙 𝑖𝑠 𝑑𝑒𝑓𝑖𝑛𝑒𝑑 𝑎𝑠
 𝓵
𝜋𝑆𝑡
𝐹𝑠 𝑓 𝑡 = 𝑓 𝑡 𝑠𝑖𝑛 𝑑𝑡; 𝑠 𝜖 𝑁
𝟎 ℓ

The above function 𝑓 𝑡 is known as inverse finite Fourier sine transform of 𝐹𝑠 𝑓 𝑡

and is defined as

∞
2 𝑠𝑖𝑛 𝜋 𝑛 𝑡
𝐹𝑆−1 𝐹𝑠 𝑓 𝑡 =𝑓 𝑡 = 𝐹𝑠 𝑓 𝑡
ℓ ℓ
𝑠=1

∞ 𝑠𝑖𝑛 𝜋 𝑛 𝑡
𝑇𝑕𝑖𝑠 𝑓𝑜𝑟𝑚𝑢𝑙𝑎 𝑖𝑠 𝑔𝑜𝑡 𝑓𝑟𝑜𝑚 𝐹𝑜𝑢𝑟𝑖𝑒𝑟 𝑠𝑖𝑛𝑒 𝑠𝑒𝑟𝑖𝑒𝑠 𝑓 𝑡 = 𝑛=1 𝑏𝑠 .
ℓ

It is also known as inversion formula.

∞ 2 𝑠𝑖𝑛 𝜋 𝑠𝑡
Hence 𝑓 𝑡 = 𝑠=1 ℓ 𝐹𝑠 𝑓 𝑡 ℓ

∞
2 𝑠𝑖𝑛 𝜋 𝑠𝑡
𝐹𝑠 𝑓 𝑡
ℓ ℓ
𝑠=1

INVERSE FUNCTION: An inverse function of a function is a function that does exactly the
opposite of the original function. If the function 𝑔 is the inverse of the function 𝑓 , and if
𝑦 = 𝑓 (𝑥), then 𝑥 = 𝑔(𝑦). It is to be noted that
1. The inverse function of a function 𝑓 ∶ 𝑋 → 𝑌 exists if and only if 𝑓 is a bijection, that
is, 𝑓 is an injection and a surjection.
2. When an inverse function exists, it is unique.
3. The inverse function and the inverse image of a set coincide in the following sense.
Suppose 𝑓 −1 (𝐴) is the inverse image of a set 𝐴 ⊂ 𝑌 under a function 𝑓 ∶ 𝑋 → 𝑌. If 𝑓 is
a bijection, then 𝑓 −1 (𝑦) = 𝑓 −1 ({𝑦}).
INVERSE FUNCTION THEOREM: For a holomorphic function 𝑓 ∶ 𝐶 → 𝐶 defined near 𝑝,
if 𝑓 ′ (𝑝) ≠ 0 then 𝑓 is a local biholomorphism near 𝑝. The theorem hands us a unique
holomorphic function 𝑔 ∶ 𝐶 → 𝐶 defined near 𝑓(𝑝) such that 𝑓(𝑔(𝑤)) = 𝑤 and
𝑔(𝑓(𝑧)) = 𝑧 (for all 𝑧, 𝑤 close enough to 𝑝, 𝑓(𝑝) respectively).
INVERSELY PROPORTIONAL: If the variables 𝑦 and 𝑥 are related by the equation
𝑦 = 𝑘/𝑥, where 𝑘 is a constant, then 𝑦 is said to be inversely proportional to 𝑥.
INVERSE MATRIX: The inverse of a square matrix 𝑨 is the square matrix of same order
that, when multiplied by 𝑨, gives the identity matrix 𝑰. Inverse of the matrix 𝑨 is written
as 𝑨−𝟏 . Note that 𝑨𝑨−𝟏 = 𝑰 = 𝑨−𝟏 𝑨. 𝑨−𝟏 exists if 𝑑𝑒𝑡 𝑨 ≠ 0.
The inverse of a 2 × 2 matrix can be found from the formula:
𝑎 𝑏 1 𝑑 −𝑏
=
𝑐 𝑑 𝑎𝑑 − 𝑏𝑐 −𝑐 𝑎
INVERSE OF A FUNCTION: Let 𝐴 and 𝐵 be sets, and let 𝑓: 𝐴 → 𝐵 be a function from 𝐴 to
𝐵. A function 𝑔: 𝐵 → 𝐴 from 𝐵 to 𝐴 is said to be the inverse of the function 𝑓 if
𝑔(𝑓(𝑎)) = 𝑎 for all elements 𝑎 of 𝐴 and 𝑓(𝑔(𝑏)) = 𝑏 for all elements 𝑏 of 𝐵. If there
exists a function 𝑔: 𝐵 → 𝐴 that is the inverse of 𝑓: 𝐴 → 𝐵, then the function 𝑓 is said to
be invertible and the inverse of a function 𝑓: 𝐴 → 𝐵 is denoted by 𝑓 −1 ∶ 𝐵 → 𝐴.

INVERSE POINTS W.R.T. A CIRCLE: Two points 𝐴(𝑧 = 𝑎) and 𝐵(𝑧 = 𝑏)are said to be
inverse points of a circle with centre 𝑂(𝑧 = 𝑧0 ) and radius 𝑟 if 𝑂. 𝐴, 𝐵 are collinear and
𝑂𝐴. 𝑂𝐵 = 𝑟 2

So that a − z0 ∙ b − z0 = r 2

Then arg a − z0 = arg b − z0 = −are b − z0

This gives arg a − z0 + arg b − z0 = 0

⇒ a − z0 b − z0 is real and equal to r 2 .

Hence a − z0 b − z0 = r 2 is the required equation.

INVERSE RULE: Suppose 𝑓 ∶ [𝑎, 𝑏] → 𝑅 is differentiable on [𝑎, 𝑏] and 𝑓 ′ (𝑥) > 0 for all
𝑥 ∇ [𝑎, 𝑏]. Let 𝑓(𝑎) = 𝑐 and 𝑓(𝑏) = 𝑑. Then the map 𝑓 ∶ [𝑎, 𝑏] → [𝑐, 𝑑] is bijective
1
and 𝑓 −1 is differentiable on [𝑐, 𝑑] with 𝑓 −1 ′ (𝑥) = .
𝑓 ′ (𝑓 −1 (𝑥))

INVERSION (COMPLEX ANALYSIS): By means of the transformation 𝑤 = 1/𝑧, figure in

𝑧 − 𝑝𝑙𝑎𝑛𝑒 are mapped upon the reciprocal figures in 𝑤 − 𝑝𝑙𝑎𝑛𝑒.

INVERTIBLE MATRICES: Let A be any n-rowed square matrix. Then a matraix 𝐵, if

exists, such that A= 𝐵𝐴 = 𝐼𝑛 is called inverse of A.

Every invertible matrix possesses a unique inverse.

IRRATIONAL NUMBER: An irrational number is a real number that is not a rational
number (i.e., it cannot be expressed as the ratio of two integers). Irrational numbers can
be represented by decimal fractions in which the digits go on forever without ever
repeating a pattern. Some of the most common irrational numbers are square roots of
positive integers, such as 3 = 1.732505 . . .. 𝜋 and 𝑒 are also irrational.
IRREDUCIBLE POLYNOMIALS: Let 𝑘 be a field. A polynomial 𝑓 𝑋 𝜖𝑘[𝑋] of degree 𝑛 is
said to be reducible over 𝑘 𝑖𝑓 𝑓 is divisible by a polynomial of degree 𝑣 < 𝑛 in
𝑘 𝑋 𝑣 ≠ 0 ; otherwise, it is said to be irreducible over 𝑘. Any polynomial of degree 1 is
irreducible. A polynomial 𝑓 is a prime element of 𝑘 𝑋 if and only if 𝑓 is irreducible over
𝑘. Let 𝐼 be a unique factorization domain. If 𝑓(𝑋) is a polynomial (1) in 𝐼 𝑋 such that for
a prime element 𝑝 in 𝐼, 𝑎𝑛 ≢ 0 (mod 𝑝), 𝑎𝑛−1 ≡ 𝑎𝑛−2 ≡ ⋯ ≡ 𝑎0 ≡ 0 mod 𝑝 but 𝑎0 ≢
0 mod 𝑝2 , then 𝑓(𝑋) is irreducible over the field of quotients of 𝐼 Eisenstein’s
theorem). If a polynomial (2) in 𝑚 variables over an algebraic number field 𝑘 is
irreducible, we can obtain an irreducible polynomial in 𝑋1 , … , 𝑋𝜇 (0 < 𝜇 < 𝑚) from the
polynomial 𝐹(𝑋1 , … , 𝑋𝑚 ) by assigning appropriate values in 𝑘 to 𝑋𝜇 +1 , … , 𝑋𝑚 .

Two lines through a vertex of a triangle are isogonal conjugate if they are
symmetric in the bisector of the triangle's angle at this vertex.

ISOGONAL CONJUGATE: Two lines through a vertex of a triangle are isogonal

conjugate if they are symmetric in the bisector of the triangle's angle at this vertex. If
three cevians in a triangle are concurrent in a point, their isogonal conjugates are also
concurrent in a point that is known as its isogonal conjugate.

ISOLATED POINT: Let 𝑋 be a topological space, let 𝑆 ⊆ 𝑋, and let 𝑥 ∇ 𝑆. The point 𝑥 is
said to be an isolated point of 𝑆 if there exists an open set 𝑈 ⊆ 𝑋 such that 𝑈 ∩ 𝑆 =
{𝑥}.
The set 𝑆 is isolated if every point in 𝑆 is an isolated point.

ISOLATED SINGULARITY: We say that a complex-valued function 𝑓 has an isolated

singularity at 𝑤 ∇ 𝐶 if 𝑓 is analytic in a punctured disc {𝑧 ∶ 0 < |𝑧 − 𝑤| < 𝑠}. If
𝑤 = ∞ we replace this by {𝑧 ∶ 𝑠 < |𝑧| < +∞}. The singularity is then removable/ a
pole/essential according to whether 𝑓 has a finite/infinite/nonexistent limit as 𝑧 → 𝑤.
We also say that a function 𝑓 is analytic apart from isolated singularities in a domain
𝐷 ⊆ 𝐶 if for every 𝑎 ∇ 𝐷 either 𝑓 is analytic at 𝑎 or 𝑎 is an isolated singularity.
ISOMETRY: A one-one mapping Φ from the metric space (𝑋, 𝜌) onto the metric space
(𝑋, 𝜌′) is called an isometry if, for any 𝑓, 𝑔 𝜖 𝑋.

𝜌 𝑓, 𝑔 = 𝜌′ Φ 𝑓 , Φ 𝑔 .

The two spaces are then said to be isometric.

As an example of isometry, take both the metric spaces to be the real number system
with the usual metric (i.e. R with 𝜌 𝑓, 𝑔 = 𝑓 − 𝑔 , 𝑓, 𝑔 𝜖 𝑅), under the mapping
𝑓 → 𝑓 + 𝑎, where 𝑎 is a fixed real number.

Another mapping which is also an isometry in this context is the one that takes any real
number 𝑓 𝑡𝑜 − 𝑓. A one-one linear transformation 𝑇 of 𝑁 into 𝑁 ∗∗ is called an isometry if
𝑇𝑓 = 𝑓 , for all 𝑓 ∇ 𝑁, when 𝑁 ∗ is conjugate space of N.

We also say that the mapping 𝑇 is an isometric isomorphism of 𝑁 into 𝑁 ∗∗ . If such an

isometric isomorphism exists we say that 𝑁 is isometrically isomorphic to 𝑁 ∗∗ .

ISOMORPHIC SPACE: Two normed linear spaces are said to be isomorphic if there exists
a one-one correspondence between them which preserves linearity and the norm.

ISOMORPHISM EXTENSION THEOREM: Given any field 𝐹, an algebraic

extension field 𝐸 of 𝐹 and an isomorphism 𝜑 mapping 𝐹 onto a field 𝐹 ′ then 𝜑 can be
extended to an isomorphism 𝜏 mapping 𝐸 onto an algebraic
extension 𝐸 ′ of 𝐹 ′ (a subfield of the algebraic closure of 𝐹 ′ ).

ISOMORPHISM OF GRAPHS: An isomorphism between two graphs (𝑉, 𝐸) and (𝑉0 , 𝐸0 ) is

a bijective function 𝜙: 𝑉 → 𝑉0 with the following property:
For any two distinct vertices 𝑎 and 𝑏 belonging to 𝑉 , {𝑎, 𝑏} ∇ 𝐸 if and only if
{𝜙(𝑎), 𝜙(𝑏)} ∇ 𝐸0 . If there exists such an isomorphism 𝜙: 𝑉 → 𝑉0 between two graphs
(𝑉, 𝐸) and (𝑉0 , 𝐸0 ) then these graphs are said to be isomorphic.
ISOPERIMETRIC PROBLEMS: It is necessary to make a given integral
𝑥
𝐼 = ∫𝑥 2 𝐼 𝑥, 𝑦, 𝑦′ 𝑑𝑥 = constant. Such problem involves one or more constraint
1

conditions. This type of problems are called isoperimetric problem.

J
JACKSON THEOREM: Let 𝑈𝑛 be a Chebyshev subspace of 𝐶[𝑎, 𝑏]. Then each 𝑓 ∇
𝐶[𝑎, 𝑏] possesses a unique polynomial of best approximation in𝐿1 -norm.
JACOBIAN: If 𝑓 (𝑥, 𝑦), 𝑔(𝑥, 𝑦) are two functions of two variables, then the Jacobian
matrix is the matrix of partial derivatives:
𝜕𝑓 𝜕𝑓
𝜕𝑥 𝜕𝑦
𝜕𝑔 𝜕𝑔
𝜕𝑥 𝜕𝑦
The equivalent definition also applies to cases with more than two dimensions. For
three functions 𝑓 𝑥, 𝑦 , 𝑔 𝑥, 𝑦 , 𝑕(𝑥, 𝑦), the Jacobian matrix becomes
𝜕𝑓 𝜕𝑓 𝜕𝑓
𝜕𝑥 𝜕𝑦 𝜕𝑧
𝜕𝑔 𝜕𝑔 𝜕𝑔
𝜕𝑥 𝜕𝑦 𝜕𝑧
𝜕𝑕 𝜕𝑕 𝜕𝑕
𝜕𝑥 𝜕𝑦 𝜕𝑧
And so on. The determinant of this matrix is known as the Jacobian determinant.
JACOBIAN OF TRANSORMATION (COMPLEX ANALYSIS): In general, the transformation

𝑢 = 𝑢 𝑥, 𝑦 , 𝑣 = 𝑣(𝑥, 𝑦)

Maps a closed region 𝑅 of z-plane upon a closed region 𝑅′ of 𝜔 − 𝑝𝑙𝑎𝑛𝑒. Let 𝛿𝐴𝑧 and
𝛿𝐴𝜔 denote respectively areas of these regions. Then it is easy to show that

𝛿𝐴𝜔
lim =𝐽
𝛿𝐴 𝑧 →0 𝛿𝐴𝑧

Provided, 𝑢, 𝑣, are continuously differentiable, where

𝜕𝑢 𝜕𝑢
𝜕(𝑢, 𝑣) 𝜕𝑥 𝜕𝑦
𝒥= =
𝜕(𝑥, 𝑦) 𝜕𝑣 𝜕𝑣
𝜕𝑥 𝜕𝑦

The determinant 𝒥 is called the Jacobian of the transformation.

JACOBI’S ELLIPTIC FUNCTIONS: The elliptic function with two simple poles and two
simple zeros in a cell are called Jacobi’s elliptic functions.

JACOBSON–BOURBAKI THEOREM: Suppose that 𝐿 is a division ring. The Jacobson–

Bourbaki theorem states that there is a natural 1: 1 correspondence between:

 Division rings 𝐾 in 𝐿 of finite index 𝑛 (in other words 𝐿 is a finite-dimensional left

vector space over 𝐾).
 Unital 𝐾-algebras of finite dimension 𝑛 (as 𝐾-vector spaces) contained in the
ring of endomorphisms of the additive group of 𝐾.

The sub division ring and the corresponding subalgebra are each other's commutants.

JENSEN'S COVERING THEOREM: Jensen's covering theorem states that if 0 does not
exist then every uncountable set of ordinals is contained in a constructible set of the
same cardinality. Informally this conclusion says that the constructible universe is close
to the universe of all sets.

JENSEN’S INEQUALITY COMPLEX ANALYSIS): Let 𝑓(𝑧) be an integral function which

does not vanish at the origin. Also let 𝑟1 , 𝑟2 , 𝑟3, … … 𝑟𝑛 be the moduli of zeros 𝑧1 , 𝑧2 , … . 𝑧𝑛
of 𝑓 𝑧 , arranged as non-decreasing sequence, multiple zero being repeated. Then

𝑅 𝑛 𝑓(𝑜) ≤ 𝑀 𝑅 𝑟1 , 𝑟2 … , 𝑟𝑛 𝑖𝑓 𝑟𝑛 < 𝑅 < 𝑟𝑛+1.

JENSEN’S FORMULA COMPLEX ANALYSIS): If 𝑓(𝑧) is analytic within and on the circle 𝛾
such that 𝑧 = 𝑅, 𝑎𝑛𝑑 𝑖𝑓 (𝑧) has zeros at the points 𝑎𝑗 ≠ 0, (𝑗 = 1,2, … . , 𝑚) and poles at
𝑏𝑗 ≠ 0, (𝑗 = 1,2, … … , 𝑛) inside 𝛾, multiple zeros and poles being repeated, then

2𝜋 𝑚 𝑛
1 𝑅 𝑅
log 𝑓 𝑅𝑒 𝑖𝜃 𝑑𝜃 = log 𝑓(0) + log − log
2𝜋 0 𝑎𝑟 𝑏𝑠
𝑟=1 𝑠=1

JENSEN’S THEOREM COMPLEX ANALYSIS): Let 𝑓(𝑧) be analytic for 𝑓(𝑧) ≤ 𝑅. Let
𝑟1 , 𝑟2 , … 𝑟𝑛 , be moduli of the zeros of 𝑓(𝑧) in 𝑧 < 𝑅 arranged as a non-decreasing
sequence. Then, 𝑖𝑓 𝑟𝑛 ≤ 𝑟 < 𝑟𝑛+1 , prove that

2𝜋
𝑟 𝑛 𝑓(0) 1
log = log 𝑓 𝑟𝑒 𝑖𝜃 𝑑𝜃
𝑟1 , 𝑟2 , … 𝑟𝑛 2𝜋 0
JOCHIMSTHAL’S THEOREM DIFFERENTIAL GEOMETRY): If the curve of intersection of
two surface is a line of curvature on both the surfaces, then the surface cut a constant
angle. Conversely, if two surfaces cut a constant angle and the curve of intersection is a
line of curvature on one of them, then it is also a line curvature on the other.

JOCKEYING: If there are number of queues, then one may leave one queue to join
another.

JOHNSON’S METHDO For n-jobs 2- machines): Let suppose 𝑛 𝑗𝑜𝑏𝑠 1,2, … 𝑛 are to be
processed on two machines say 𝐴 𝑎𝑛𝑑 𝐵 𝑎𝑛𝑑 𝐴𝑖 , 𝐵𝑖 , 𝑖 = 1,2, … 𝑛 are the respective
processing times of 𝑖 𝑡𝑕 job on 𝐴 𝑎𝑛𝑑 𝐵 machines respectively.

Assumptions

(i) Each job is processed in order AB.

(ii) 𝐴𝑖 = Processing time of 𝑖 𝑡𝑕 job on machine 𝐴 𝑖 = 1,2 … 𝑛
(iii) 𝐵𝑖 = Processing time of 𝑖 𝑡𝑕 job on machine 𝐵 𝑖 = 1,2 … 𝑛

We want to find the sequence of jobs to be performed on two machines so that the total
time 𝑇 elapsed from the start of the first job to the completion of the last job to be
minimized.

Step1. Select the smallest processing time in the list 𝐴1 , 𝐴2 , … 𝐴𝑛 𝑎𝑛𝑑 𝐵1 , 𝐵2 , … 𝐵𝑛 . If there
is a tie then either of these smallest processing time may be selected or in this case
consider the following cases:

(i) Minimum of all the processing times is 𝐴𝑟 which is also equal to 𝐵𝑠 . Then min
𝐴𝑖 , 𝐵𝑖 = 𝐴𝑟 = 𝐵𝑠 . . Then do the 𝑟 𝑡𝑕 job first and 𝑠 𝑡𝑕 job in the end.
(ii) If min 𝐴𝑖 , 𝐵𝑖 = 𝐴𝑟 but also 𝐴𝑟 = 𝐴𝑘 𝑠𝑎𝑦 then do anyone of these jobs for
which there is a tie, first.
(iii) If there is a tie for minimum among 𝐵𝑖′ 𝑠𝑖. 𝑒. , 𝑀𝑖𝑛 𝐴𝑖 , 𝐵𝑖 = 𝐵𝑠 = 𝐵𝑟 𝑠𝑎𝑦 then
do anyone of these jobs in the last.

Step2. If the smallest processing time is 𝐴𝑟 𝑖. 𝑒. , 𝑖𝑛 𝑡𝑕𝑒 𝑙𝑖𝑠𝑡 𝐴1 , … , 𝐴𝑛 then do the 𝑟 𝑡𝑕

job first. On the other hand if it is 𝐵𝑠 𝑖. 𝑒. , 𝑖𝑛 𝑡𝑕𝑒 𝑙𝑖𝑠𝑡 𝐵1 , 𝐵2 , … , 𝐵𝑛 . Then do the 𝑠 𝑡𝑕 job
last.
Step 3. Delete the already assigned job from both the list. If 𝑦 𝑡𝑕 job is assigned
previously, then delete 𝐴𝑟 𝑎𝑛𝑑 𝐵𝑟 both and if 𝑠 𝑡𝑕 job is assigned previously hen delete
𝐴𝑠 𝑎𝑛𝑑 𝐵𝑠 both.

Step4. Repeat step (1) to (3) for remaining jobs.

Step5. Continuing the same process until all the jobs have been ordered and get optimal
sequence of jobs.

JOINT VARIATION: If 𝑧 = 𝑘𝑥𝑦, where k is a constant, then 𝑧 is said to vary jointly with x
and y.
JORDAN ARC: Suppose 𝑥 t and 𝑦 t are continuous real valued functions in the range
𝛼 ≤ 𝑡 ≤ 𝛽, then the set of points 𝑧 in the Argand plane determined by the equation
𝑧 = 𝑥 t + i𝑦 t is called a continuous arc. A point z1 is called multiple point of the arc if
the equation 𝑧 = 𝑥 t + i𝑦 t is called a continuous arc. A point z1 is called multiple
point of the arc if the equation 𝑧1 = 𝑥 t + i𝑦 t is satisfied by more than one value of t
in the given range. A multiple point z1 of the arc is called a double point if the equation
(1) is satisfied by two values of t in the given range. A continuous arc without multiple
points is called a Jordan arc.

If the end points t = α and t = β of the continuous arc are coincident and if the arc has
only one multiple point (which is the double point corresponding to the two terminal
values of t), then the arc is called simple closed Jordan closed curve.

JORDAN-BROUWER THEOREM: Let 𝑀 ⊂ 𝑅 𝑛 be an 𝑛 − 1 –dimensional compact

connected manifold in 𝑅 𝑛 . The complement of 𝑀 in 𝑅 𝑛 consists of precisely two
components, of which one, called the outside is unbounded, and the other, called the
inside is bounded. Each of the two components is a domain with smooth boundary 𝑀.
JORDAN CURVE THEOREM: It states that a simple closed Jordan curve divides the
Argand plane into two open domains which have the curve as the common boundary.

One of these domain is bounded and is called interior domain while the other is
unbounded and is called exterior domain.
JORDAN DECOMPOSITION THEOREM: Let ν be an additive set function on an algebra 𝐴
of subsets of 𝑋, and suppose that 𝜈 has finite total variation. Then, for all 𝐴 ∇ 𝐴, 𝜈(𝐴) =
𝑉 (𝜈, 𝐴) + 𝑉 (𝜈, 𝐴) .
JORDAN–HÖLDER THEOREM: Any two composition series of a given group are
equivalent. That is, they have the same composition length and the same composition
factors, up to permutation and isomorphism.

JORDAN’S INEQUALITY: Jordan’s Inequality states that

2 𝜋
𝑥 ≤ 𝑆𝑖𝑛 𝑥 ≤ 𝑥 ∀ 𝑥 ∇ 0,
𝜋 2
JORDAN’S LEMMA: If 𝑓(𝑧) is analytic except at finite number of singularities and if
𝑓 𝑧 ⟶ 0 unformaly as 𝑧 ⟶ ∞, then

lim ∫Γ 𝑒 𝑖𝑚𝑧 𝑓 𝑧 𝑑𝑧 = 0, 𝑚 > 0

𝑅→∞

where Γ denotes the semi-circle 𝑧 = 𝑅, 𝐼 𝑧 > 0

Here 𝑅 is taken so large that all the singularities of 𝑓(𝑧) lie within the semi circle Γ. (No
singularity lies on the boundary of the semi circle).

JORDAN NORMAL FORM: Let 𝑉 be a finite dimensional vector space over a complete
field 𝐹 and α an endomorphism of 𝑉 . Then there is basis for 𝑉 such that 𝛼 has matrix 𝐴
(relative to this basis) which consists of zeros except for Jordan matrices 𝐽(𝜆𝑖 , 𝑛𝑖 ) [1 ≤
𝑖 ≤ 𝑠] down the diagonal.

JORDAN NORMAL FORM FOR MATRICES: If 𝐴 is a 𝑛 × 𝑛 matrix over a complete field

then we can find an invertible 𝑛 × 𝑛 matrix 𝑃 such that 𝐽𝐴 = 𝑃𝐴𝑃 −1 consists of zeros
except for Jordan matrices 𝐽(𝜆𝑖 , 𝑛𝑖 ) [1 ≤ 𝑖 ≤ 𝑠] down the diagonal. The matrix 𝐽𝐴 so
associated with 𝐴 is unique up to reordering the diagonal blocks.

c2
JOUKOWSKI TRANSFORMATION: The transformation ζ = z + is called the Joukowski
z

transformation mapping function from 𝑧- plane on the 𝜁- plane.

JULIA SET: The set of points for a function of the form 𝑧 2 + 𝑐 (where 𝑐 is a complex
parameter), such that a small perturbation can cause drastic changes in the sequence of
iterated function values and iterations will either approach zero, approach infinity or
get trapped in loop.
K

KELVIN’S CIRCULATION THEOREM: Kelvin’s theorem states that the circulation round
any closed material curve is invariant in an inviscid fluid provided that the body force
are conservative and the pressure is single-valued function of density only.

KELVIN’S MINIMUM ENERGY THEOREM: The irrotational motion of a liquid occupying a

simply connected region has less kinetic energy than any other motion consistent with
the same normal velocity of the boundary.

KENDALL’S RANK CORRELATION: Take pairs 𝑅𝑖 , 𝑆𝑖 and 𝑅𝑗 , 𝑆𝑗 . If 𝑅𝑗 − 𝑅𝑖 𝑆𝑗 − 𝑆𝑖 >

𝑛 −1
0, set 𝜑𝑖𝑗 = 1; otherwise, 𝜑𝑖𝑗 = −1. This statistic 𝑟𝑘 = 2
𝜑𝑖𝑗 is called Kendall’s
correlation, where runs over all possible pairs chosen from 𝑅1 , 𝑆1 , … , 𝑅𝑛 , 𝑆𝑛 .

If there is no dependence between 𝑋 and 𝑌, then 𝑟𝑘 = 0 and

𝑉 𝑟𝑘 = 2 2𝑛 + 5 9𝑛 𝑛 − 1 .

KEPLER, JOHANNES: Johannes Kepler (1571-1630) was a German astronomer who used
observational data to express the motion of the planets according to three mathematical
laws:
 Planets move along orbits shaped like ellipses, with the sun at one focus;
 A radius vector connecting the sun to the planet sweeps out equal areas in equal
times (this means that a planet travels fastest when closest to the sun);
 The square of the orbital period is proportional to the cube of the mean distance
from the planet to the sun.
KERNEL OF LINEAR FUNCTIONAL: The kernel of linear functional 𝛼, write 𝑘𝑒𝑟𝛼, is the
set all vectors 𝑥 ∇ 𝑋 such that 𝛼(𝑥) = 0.

Note that

1. X* is a linear space with natural operations.

2. 𝛼 𝑥 ≤ 𝛼 · 𝑥 𝑓𝑜𝑟 𝑎𝑙𝑙 𝑥 ∇ 𝑋, 𝛼 ∇ 𝑋 ∗.

3. X* is a Banach space with the defined norm (even if X was incomplete).

4. 𝑘𝑒𝑟𝛼 is a subspace of 𝑋.

5. If 𝛼 ≢ 0 then 𝑘𝑒𝑟𝛼 is a proper subspace of 𝑋.

 6. If 𝛼 is continuous then 𝑘𝑒𝑟𝛼 is closed.
1
KINETIC ENERGY: We define the kinetic energy 𝑇 as 𝑇 = 𝑚 𝑟 . 𝑟 , where 𝑟 is the
2

velocity of the particle.

KIRCHHOFF VORTEX THEOREM: Kirchhoff vortex theorem states that if

r1 , θ1 , r2 , θ2 , … … are the polar coordinates of number of rectilinear vortices of
strengths k1 , k 2 , …. then 𝑛 k i , k i = A1 , k i yi = A2 , k i ri2 = A3 , k i ri2 θi = A4 , where
A1 , A2 , A3 , A4 are constants.

KNOT THEORY: An area of topology that studies mathematical knots (a knot is a closed
curve in space formed by interlacing a piece of “string” and joining the ends .

KOLMOGOROV THEOREM: Let 𝑈 be a linear subspace of 𝐶(𝐾). An element 𝑝∗ ∇ 𝑈 is a

best approximation to 𝑓 ∇ 𝐶(𝐾) if and only if max 𝑥 ∇ 𝑍 [𝑓(𝑥) − 𝑝∗ (𝑥)] 𝑞(𝑥) ≥
0 ∀𝑞 ∇ 𝑈 , where 𝑍 is the set of all points for which |𝑓(𝑥) − 𝑝∗ (𝑥)| = 𝑘𝑓 − 𝑝∗ 𝑘.

KÖNIG'S THEOREM: In any bipartite graph, the number of edges in a maximum

matching equals the number of vertices in a minimum vertex cover.

KRONECKER THEOREM: Let 𝐾 be a field, and let 𝑓 ∇ 𝐾[𝑥] be a nonconstant polynomial

with coefficients in 𝐾. Then there exists an extension field 𝐿 of 𝐾 and an element 𝛼 of 𝐿
for which 𝑓(𝛼) = 0.

KRULL'S PRINCIPAL IDEAL THEOREM: It gives a bound on the height of a principal

ideal in a Noetherian ring. Formally, if 𝑅 is a Noetherian ring and 𝐼 is a principal, proper
ideal of 𝑅, then 𝐼 has height at most one.

This theorem can be generalized to ideals that are not principal, and the result is often
called Krull's height theorem. This says that if 𝑅 is a Noetherian ring and 𝐼 is a proper
ideal generated by 𝑛 elements of 𝑅, then 𝐼 has height at most 𝑛.

KUHN- TUCKER CONSTRAINT QUALIFICATION: Let X 0 be an open set in r n and let 𝑔 be

an 𝑚- dimensional vector function defined on X 0 . Let

𝑋 = x: x ∇ X 0 , g(x) ≤ 0
𝑔 is said to satisfy the Kuhn- Tucker constraint qualification at x ∇ X if 𝑔 is differentiable
at x and if there exist an n – dimensional vector function e defined on the interval [0,1]
such that

e 0 = x, e τ ∇ X for 0 ≤ τ ≤ I

𝑑𝑒 (0)
and e is differentiable at 𝜏 = 0 and = 𝜆𝛾 for some 𝜆 > 0 where 𝛾 ∇ 𝑅 𝑛 ,
𝑑𝜏

∆g 𝑙 (x)𝛾 ≤ 0 and I = i g i x = 0

KUHN- TUCKER STATIONARY- POINT PROBLEM (KTP): The Kuhn- Tucker stationary-
point problem means to find x ∇ X 0 , u ∇ Rm if they exist, such that ∆x Ψ x, u = 0,
∆u Ψ x, u ≤ 0, u∆u Ψ x, u = 0, u ≥ 0 and Ψ x, u = θ x) + ug(x or equivalently,
∆θ x) + u ∆g(x = 0, g(x) ≤ 0, u g x = 0 and u ≥ 0

KURT GÖDEL: Kurt Gödel was born in the Czech Republic and grew up in Austria (which
was the Austro-Hungarian Empire in Gödel's early childhood). His primary language
was German. Although he is most famous for his contribution to mathematical logic, he
also did much work in set theory . He was a friend of Albert Einstein during the time
they were both at the Institute for Advanced Study at Princeton University

KUTTA- JOUKOWSKI’S THEOREM: When an aerofoil at rest be placed in a uniform wind

of speed 𝑈 with circulation 𝑘 around the aerofoil, if experience a lift of magnitude 𝑝𝑘𝑈
per span perpendicular to the wind. The direction of the lift is obtained by rotating the
wind velocity through a right angle in the sense opposite to that of the circulation.

L
LAGRANGE, JOSEPH-LOUIS: Joseph-Louis Lagrange (1736-1813) was an Italian-French
mathematician who developed ideas in celestial mechanics, calculus of variations, and
number theory.

LAGRANGE'S FOUR-SQUARE THEOREM: Lagrange's four-square theorem, also known

as Bachet's conjecture, states that any natural number can be represented as the sum of
four integer squares.

where the four numbers are integers. For illustration, 3, 31 and 310 can
be represented as the sum of four squares as follows:

LAGRANGE’S EQUATION: Consider a conservative holomonic dynamical system defined

by n generalized coordinates 𝑞1 , 𝑞2 , … . . 𝑞𝑛 at time t. The Lagrangian function L=T-V
where T and V are kinetic and potential energies respectively of the system is a function
of 𝑞1 , 𝑞2 , … . . 𝑞𝑛 𝑎𝑛𝑑 𝑞1 , 𝑞2 , … . . 𝑞𝑛 at time t.
𝑡
By Hamiltonian’s principal ∫𝑡 1 𝐿 𝑑𝑡 is stationary.
0

Therefore, Euler’s equation must hold good.

i.e.,
𝜕𝐿 𝑑 𝜕𝐿
− = 0, 𝑖 = 1,2 … , 𝑛.
𝜕𝑞𝑖 𝑑𝑡 𝜕𝑞𝑖
These are called Lagrange’s equation which determine the motion of the system.
Proper field: A family of curves 𝑦 = 𝑦 𝑥, 𝑐 is said to form a proper field for a domain D
of the 𝑥𝑦- plane if through any point of D, there passes one and only one curve of the
family.
Consider the circle 𝑥 2 + 𝑦 2 = 1 . Inside the circle 𝑥 2 + 𝑦 2 ≤ 1, the family of curve
𝑦 = 𝑚𝑥 + 𝑐 is a proper field, because there passes one and only one curve of the family
through any point of the given circle.
LAGRANGE’S MEAN VALUE THEOREM: If a real valued function 𝑓(𝑥) is such that
 𝑓(𝑥) is continuous in the closed interval [𝑎, 𝑏]
 𝑓(𝑥) is differentiable in the open interval (𝑎, 𝑏)
 𝑓 𝑏 −𝑓(𝑎)
then there exist at least one value of 𝑥 = 𝑐 ∇ (𝑎, 𝑏) such that 𝑓 ′ (𝑐) = .
𝑏−𝑎

LAGRANGE MULTIPLIERS: The method of Lagrange multipliers is used to find maximum

or minimum values in the presence of constraints. For example, suppose we need to
choose x and y to maximize the function
𝑧 𝑥, 𝑦 = 𝑎𝑥 + 𝑏𝑦
subject to this constraint:
𝑅 − 𝑝𝑥 2 − 𝑞𝑦 2 = 0
Create the Lagrangian function 𝐿 as follows:
𝐿 = 𝑎𝑥 + 𝑏𝑦 + 𝜇(𝑅 − 𝑝𝑥 2 − 𝑞𝑦 2 )
Notice that the first part of the Lagrangian is the function we are trying to maximize.
The second part consists of a new variable 𝜇 (called the Lagrange multiplier), multiplied
by the left-hand side of the constraint equation. (To do this, arrange the equation so that
the right-hand side is zero.) The method also works with more than one constraint; just
add a new Lagrange multiplier for each one. Now find the partial derivatives of 𝐿 with
respect to 𝑥, 𝑦 and 𝜇 , and set them all equal to zero:
𝜕𝐿
= 𝑎 − 2𝜇𝑝𝑥 = 0
𝜕𝑥
𝜕𝐿
= 𝑏 − 2𝜇𝑞𝑦 = 0
𝜕𝑦
𝜕𝐿
= 𝑅 − 𝑝𝑥 2 − 𝑞𝑦 2 = 0
𝜕𝜇
Now solve this three-equation system.
LAGRANGIAN OF A SYSTEM: Let 𝑇 be the kinetic energy and 𝑉 be the potential energy of
a system then, the expression 𝐿 = 𝑇 − 𝑉 is known as Lagrangian of a system.
LAGRANGIAN SUFFICIENCY THEOREM: Suppose we are given a general optimization
problem, P:

Minimize 𝑓(𝑥) s.t. 𝑔(𝑥) = 𝑏, 𝑥 ∇ 𝑋, with 𝑥 ∇ 𝑅 𝑛 , 𝑏 ∇ 𝑅 𝑚 (𝑛 variables and 𝑚

constraints).

The Lagrangian is 𝐿(𝑥, 𝜆) = 𝑓(𝑥) − 𝜆 ⊤ (𝑔(𝑥) − 𝑏), with 𝜆 ∇ 𝑅 𝑚 (one component for
each constraint).
Each component of 𝜆 is called a Lagrange multiplier. The following theorem is simple to
prove, and extremely useful in practice. If 𝑥 ∗ and 𝜆∗ exist such that 𝑥 ∗ is feasible for
𝑃 and 𝐿(𝑥 ∗ , 𝜆∗ ) ≤ 𝐿(𝑥, 𝜆∗ ) ∀𝑥 ∇ 𝑋, then 𝑥 ∗ is optimal for 𝑃.

LAGUERRE DIFFERENTIAL EQUATION: The differential equation

d2 y dy
𝑥 2 + 1−𝑥 + λy = 0,
dx dx

where λ is a constant, is called Laguerre differential equation.

LAGUERRE POLYNOMIALS: We define the standard solution of Laguerre equation

d2 y dy
𝑥 2 + 1−𝑥 + ny = 0
dx dx

As that for which a0 = 1 and call it the Laguerre polynomial of order 𝑛, denote it by

𝑛
𝑛!
𝐿𝑛 𝑥 , = (−1)𝑟 𝑥′
𝑛 − 𝑟 ! (𝑟!)2
𝑟=0

LAMBERT SERIES: A series of the form

∞ 𝑧𝑛
𝑛=1 𝑎𝑛 1−𝑧 𝑛 , 𝑧 ≠ 1, (1)

is called a Lambert series. If 𝑎𝑛 is convergent, (1) converges for any 𝑧 with 𝑧 ≠ 1,

and moreover, it converges uniformly on any compact ser contained in 𝑧 < 1 or 𝑧 >
1. If 𝑎𝑛 is divergent, (1) and the power series 𝑎𝑛 𝑧 𝑛 converge or diverge
simultaneously for 𝑧 ≠ 1.

LAMI'S THEOREM: Lami's theorem is an equation relating the magnitudes of

three coplanar, concurrent and non-collinear forces, which keeps an object in static
equilibrium, with the angles directly opposite to the corresponding forces. According to
the theorem,
 where 𝐴, 𝐵 and 𝐶 are the magnitudes of three coplanar, concurrent and non-collinear
forces, which keep the object in static equilibrium, and 𝛼, 𝛽 and 𝛾 are the angles
directly opposite to the forces 𝐴, 𝐵 and 𝐶 respectively.
LANDAU’S THEOREM: Let 𝐸 = {𝑧 ∶ |𝑧| ≤ 1}. Suppose that 𝑓 is analytic on 𝐸, and that
𝑓 ′ (0) = 1. Then 𝑓(∆) contains a disc of radius 1/24.
LANGUAGE: Let 𝐴 be a finite set. A language over 𝐴 is a subset of 𝐴∗ . A language 𝐿 over 𝐴
is said to be a formal language if there is some finite set of rules or algorithm that will
generate all the words that belong to 𝐿 and no others.
LAPLACE, PIERRE-SIMON: Pierre-Simon Laplace (1749- 1827) was a French astronomer
and mathematician who investigated the motion of the planets of the solar system.
LAPLACE’S FIRST INTEGRAL FOR 𝐏𝐧 (𝐱): If n is a positive integer, then

1 π n
Pn x =π ∫0 x ± x 2 − 1 cos∅ d∅.

LAPLACE’S SECOND INTEGRAL FOR 𝐏𝐧 (𝐱): If n is a positive integer, then

1 π d∅
Pn x =π ∫0 n +1 .
x± x 2 −1 cos∅

LAPLACE TRANSFORMATION:
Laplace
Entry Function
Tranformation
unit impulse unit impulse

unit step

slope

parabola

tn
(n is integer)

exponential

time
multiplied
exponential
Asymptotic
exponential
double
exponential
asymptotic
double
exponential
asymptotic
critically
damped
differentiated
critically
damped

sine

cosine

decaying
sine
decaying
cosine
generic
decaying
oscillatory

generic
decaying
oscillatory
(alternate)
Prototype
2nd order
lowpass
step
response
Prototype
2nd order
lowpass
impulse
response
Prototype
2nd order
bandpass
impulse
response
LAPLACIAN: The Laplacian of a function 𝑓 (𝑥, 𝑦, 𝑧) is:

2
∂2 f ∂2 f ∂2 f
∆ = 2+ 2+ 2
∂x ∂y ∂z
It is the divergence of the gradient of 𝑓 .
LATERAL AREA: The lateral area of a solid is the area of its faces other than its bases.
For example, the lateral area of a pyramid is the total area of the triangles forming the
sides of the pyramid.
LATUS RECTUM: The latus rectum of a parabola is the chord through the focus
perpendicular to the axis of symmetry. The latus rectum of an ellipse is one of the
chords through a focus that is perpendicular to the major axis.
LAURENT’S THEOREM COMPLEX ANALYSIS): Suppose a function 𝑓(𝑧) is analytic in the
closed ring bounded by two concentric 𝐶 𝑎𝑛𝑑 𝐶′ of centre 𝑎 and radii 𝑅 𝑎𝑛𝑑 𝑅 ′ , 𝑅 ′ <
𝑅 . If 𝑧 is any point of the annulus, then

∞ ∞

𝑓 𝑧 = 𝑎𝑛 (𝑧 − 𝑎)𝑛 + 𝑏𝑛 (𝑧 − 𝑎)−𝑛
𝑛=0 𝑛=1

Where

1 𝑓 𝑡 𝑑𝑡 1 𝑓 𝑡 𝑑𝑡
𝑎𝑛 = ∫𝐶 , 𝑏𝑛 = ∫ ′
2𝜋𝑖 (𝑡 − 𝑎)𝑛+1 2𝜋𝑖 𝐶 (𝑡 − 𝑎)−𝑛+1
LAW OF AVERAGES: The law of averages is an erroneous generalization of the law of
large numbers, which states that the frequencies of events with the same likelihood of
occurrence even out, given enough trials or instances. The law of averages is usually
mentioned in reference to situations without enough outcomes to bring the law of large
numbers into effect.

A common example of how the law of averages can mislead involves the tossing of a fair
coin (a coin equally likely to come up heads or tails on any given toss). If someone
tosses a fair coin and gets several heads in a row, that person might think that the next
toss is more likely to come up tails than heads in order to "even things out." But the true
probabilities of the two outcomes are still equal for the next coin toss and any coin toss
that might follow. Past results have no effect whatsoever: Each toss is an independent
event. The law of large numbers is often confused with the law of averages, and many
texts use the two terms interchangeably. However, the law of averages, strictly defined,

LAW OF EXCLUDED MIDDLE: One of the two: either a statement is true or it is false.
There's nothing in-between. An axiom of 2-valued logic.

LAW OF LARGE NUMBERS: The law of large numbers states that if a random variable is
observed many times, the average of these observations will tend toward the expected
value (mean) of that random variable. For example, if you roll a die many times and
calculate the average value for all of the rolls, you will find that the average value will
tend to approach 3.5.
LEADER OF A COSET (CODING THEORY): The leader of a coset is defined to be the word
with the smallest hamming weight in the coset.
LEBESGUE'S DOMINATED CONVERGENCE THEOREM: Let {𝑓𝑛 } be a sequence of real-
valued measurable functions on a measure space (𝑆, 𝛴, 𝜇). Suppose that the sequence
converges pointwise to a function 𝑓 and is dominated by some integrable function 𝑔 in
the sense that

for all numbers 𝑛 in the index set of the sequence and all points 𝑥 ∇ 𝑆. Then 𝑓 is
integrable and

which also implies

LEBESGUE INEQUALITY: Let X be a normed linear space, and let 𝑀 ∶ 𝑋 → 𝑈 be a linear

projection. Then, for every 𝑓 ∇ 𝑋, the approximation 𝑀(𝑓) has the property

𝑓 − 𝑀(𝑓) ≤ ( 𝑀 + 1) 𝑑𝑖𝑠𝑡(𝑓, 𝑈).

LEBESGUE INTEGRAL: For an arbitrary summable 𝑓 ∇ 𝐿1 (𝑋), we define the Lebesgue

integral

𝑓𝑑𝜇 = lim 𝑓𝑛 𝑑𝜇,

𝑛

where the Cauchy sequence 𝑓𝑛 of summable simple functions converges to 𝑓 a.e.

1. 𝐿1 (𝑋) is a linear space.

2. For any set 𝐴 ⊂ 𝑋 the correspondence 𝑓 ↦ ∫ 𝑓 𝑑µ is a linear functional on

𝐿1 (𝑋).

3. 𝑑1 (𝑓, 𝑔) = ∫ | 𝑓 − 𝑔 | 𝑑µ is a distance on 𝐿1 (𝑋).

LEBESGUE MEASURABLE SET: Let 𝑆 ⊆ 𝑅, and let 𝑆 ′ be the complement of 𝑆 with

respect to 𝑅. We define 𝑆 to be measurable if, for any 𝐴 ⊆ 𝑅,
𝑚∗ (𝐴) = 𝑚∗ (𝐴 ∩ 𝑆) + 𝑚∗ (𝐴 ∩ 𝑆 ′ )
where 𝑚∗ (𝑆) is the Lebesgue outer measure of 𝑆. If 𝑆 is measurable, then we define the
Lebesgue measure of 𝑆 to be 𝑚(𝑆) = 𝑚∗ (𝑆).
LEBESGUE ON DOMINATED CONVERGENCE: Let (𝑓𝑛 ) be a sequence of µ-summable
functions on 𝑋, and there is 𝜑 ∇ 𝐿1 (𝑋). such that | 𝑓𝑛 (𝑥) | ≤ 𝜑(𝑥) for all 𝑥 ∇ 𝑋, 𝑛 ∇ ℕ. If
𝑎. 𝑒
𝑓𝑛 𝑓, then 𝑓 ∇ 𝐿1 (𝑋) and for any measurable A: limn→∞ ∫ 𝑓𝑛 𝑑𝜇 = ∫ 𝑓𝑑𝜇
→
If 𝑔 is measurable and bounded then 𝑓 = 𝜑 𝑔 is µ-summable and for any µ-measurable
set A we have ∫ 𝑓𝑑𝜇 = 𝜑 ∫ 𝑔𝑑𝜇.

LEFT NULLITY OF A MATRIX: Suppose 𝑋 is an 𝑚-vector written in the form of a row

vector. Then the matrix product 𝑋𝐴 is defined. The subspaces 𝑆 or 𝑉𝑚 generated by the
row vector 𝑋 belonging to 𝑉𝑚 such that 𝑋𝐴 = 0 is called the row null space of the matrix
𝐴. The dimension 𝑠 of 𝑆 is called the left nullity or row nullity of the matrix 𝐴.

LEGENDRE DUPLICATION FORMULA: Legendre duplication formula is 𝑇 𝑚 𝑇 𝑚 + 12 =

𝜋
𝑇(2𝑚), where 𝑚 is an integer.
22𝑚 −1

LEGENDRE EQUATION: The differential equation of the form

d2 y dy
(1 − x 2 ) 2
− 2x + n n + 1 y = 0
dx dx

is called Legendre’s differential equation or Legendre’s equation , where n is a

constant. This equation can also be written as

𝑑 dy
1 − x2 + n n + 1 y = 0.
𝑑𝜇 dx

LEGENDRE EQUATION FROM LAPLACE’S EQUATION: Laplace’s equation in spherical

coordinates is

∂ 𝜕𝑉 1 𝜕 ∂V 1 ∂2 V
𝑟 2 𝜕𝑟 + 𝑠𝑖𝑛 𝜃 𝜕𝜃 sin 𝜃 ∂θ + 𝑠𝑖𝑛 2 𝜃 ∂∅2 = 0 ……….. i
∂r

Putting V = r n Un , where Un is a function of θ and ∅ only

∂V ∂V ∂U n ∂ 2 V ∂2 U n
So that ∂r = nr n−1 Un , ∂θ = r n , ∂∅2 = r n .
∂θ ∂∅2

Substituting in (i), we have

∂ 1 𝜕 ∂U n 1 ∂2 U n
𝑛𝑟 𝑛+1 𝑈𝑛 + 𝑠𝑖𝑛 𝜃 𝜕𝜃 r 2 sin 𝜃 + 𝑠𝑖𝑛 2 𝜃 . 𝑟 𝑛 =0
∂r ∂θ ∂∅2
 𝑟𝑛 𝜕 ∂U n 𝑟𝑛 ∂2 U n
Or n 𝑛 + 1 𝑟 𝑛 𝑈𝑛 + 𝑠𝑖𝑛 𝜃 𝜕𝜃 sin 𝜃 + 𝑠𝑖𝑛 2 𝜃 . 𝑟 𝑛 =0
∂θ ∂∅2

1 𝜕 ∂U n
Or n 𝑛 + 1 𝑈𝑛 + 𝑠𝑖𝑛 𝜃 𝜕𝜃 sin 𝜃 = 0. ………….. ii
∂θ

Supposing 𝑈𝑛 to be independent of ∅.

Putting μ = cos θ and Un = y

∂U n ∂y ∂ y ∂μ ∂y
So that = ∂θ = . = − sin θ ∂μ .
∂θ ∂μ ∂θ

Substituting in (ii), we have

1 𝜕 ∂y
n 𝑛 + 1 𝑦 + 𝑠𝑖𝑛 𝜃 𝜕𝜃 −sin2 𝜃 ∂μ = 0

1 𝜕 ∂y ∂μ
or n 𝑛 + 1 𝑦 + 𝑠𝑖𝑛 𝜃 𝜕𝜇 − 1 − μ2 . ∂θ = 0
∂μ

𝜕 ∂y
or − 1 − μ2 + n n + 1 y = 0.
𝜕𝜇 ∂μ

which is Legendre’s equation.

Solution Pn (μ) of Legendre’s equation is the surface spheric harmonic of degree 𝑛, which
is free from ∅.

LEGENDRE FUNCTION OF THE SECOND KIND: We define the Legendre’s function of the
second kind, Qn , such that

n! n + 1 (n + 2) −n−3
Qn x = x −n−1 + ∙x
1 ∙ 3 … (2n + 1) 2 ∙ (2n + 3)
n + 1 n + 2 n + 3 (n + 4) −n−5
+ x
2 ∙ 4 2n + 3 (2n + 5)

2n n! ∞ (n+2r)x −(n +2r +1)

or Qn x = (2n+1)! r=0 r! 2n+3 2n+5 … 2n+2r+1 .

LEGENDRE SYMBOL: Fermat proved that if p is a prime number and a is an integer,

Thus, if p does not divide a,

Legendre lets 𝑎 and 𝐴 represent positive primes congruent to 1 (mod 4)
and 𝑏 and 𝐵 positive primes congruent to 3 (mod 4), and sets out a table of eight
theorems that together are equivalent to quadratic reciprocity:

Theorem When it follows that

II

III

IV

VI

VII

VIII

He says that since expressions of the form

(where 𝑁 and 𝑐 are relatively prime) will come up so often he will abbreviate them as:

This is now known as the Legendre symbol, and an equivalent definition is used today:

For all integers 𝑎 and all odd primes 𝑝

LEIBNIZ: Gottfried Wilhelm Leibniz (1646 to1716) was a German mathematician,

philosopher and political advisor, who was one of the developers of calculus.
LEMMA: A lemma is a theorem that is proved mainly as an aid in proving another
theorem. There is no technical distinction between a lemma and a theorem. A lemma is
a proven statement, typically named a lemma to distinguish it as a truth used as a
stepping stone to a larger result rather than an important statement in and of itself. Of
course, some of the most powerful statements in mathematics are known as lemmas,
including Zorn’s lemma, Bezout’s lemma, Gauss’ lemma, Fatou’s lemma, etc., so one
clearly can’t get too much simply by reading into a proposition’s name.
LEMNISCATES: A lemniscate is a plane curve with a characteristic shape, consisting of
two loops that meet at a central point as shown below. The curve is also known as the
lemniscate of Bernoulli.

In the (x, y) plane, the lemniscate can be described in terms of the following general
equation:

(𝑥 2 + 𝑦 2 ) 2 = 2𝑎2 (𝑥 2 − 𝑦 2 )

where 𝑎 represents the greatest distance between the curve and the origin. There are
two points on the curve that meet this criterion; both of them lie on the x axis. In the
above graph, if each division represents one unit, then 𝑎 = 5.

The lemniscate, reduced in size to that of typographical characters, is commonly used as

the symbol for infinity, or for a value that increases without limit.

LIE ALGEBRA: A Lie algebra is a vector space g over R, equipped with a map
[·,·]: 𝑔 × 𝑔 → 𝑔 satisfying the properties
 [𝑎𝑋 + 𝑏𝑌, 𝑍] = 𝑎[𝑋, 𝑍] + 𝑏[𝑌, 𝑍], [𝑋, 𝑎𝑌 + 𝑏𝑍] = 𝑎[𝑋, 𝑌 ] + 𝑏[𝑋, 𝑍],
 [𝑋, 𝑌 ] = −[𝑌, 𝑋],
 [[𝑋, 𝑌 ], 𝑍] + [[𝑌, 𝑍], 𝑋] + [[𝑍, 𝑋], 𝑌 ] = 0,
for all 𝑋, 𝑌, 𝑍 ∇ 𝑔 and 𝑎, 𝑏 ∇ 𝑅.
LIE BRACKET: Let 𝑋, 𝑌 ∇ 𝑋(𝑀) be smooth vector fields on 𝑀. The Lie bracket [𝑋, 𝑌 ] is
the linear operator 𝐶 ∞ (Ω) → 𝐶 ∞ (Ω) defined by the equation
[𝑋, 𝑌 ]𝑓 = 𝑋𝑌 𝑓 − 𝑌 𝑋𝑓 𝑓𝑜𝑟 𝑓 ∇ 𝐶 ∞ (Ω).
The Lie bracket satisfies
 [𝑎𝑋 + 𝑏𝑌, 𝑍] = 𝑎[𝑋, 𝑍] + 𝑏[𝑌, 𝑍], [𝑋, 𝑎𝑌 + 𝑏𝑍] = 𝑎[𝑋, 𝑌 ] + 𝑏[𝑋, 𝑍],
 [𝑋, 𝑌 ] = −[𝑌, 𝑋],
 [[𝑋, 𝑌 ], 𝑍] + [[𝑌, 𝑍], 𝑋] + [[𝑍, 𝑋], 𝑌 ] = 0,
 [𝑓𝑋, 𝑔𝑌 ] = 𝑓𝑔[𝑋, 𝑌 ] + 𝑓(𝑋𝑔)𝑌 − 𝑔(𝑌 𝑓)𝑋.
for all 𝑋, 𝑌, 𝑍 ∇ 𝑋(𝑀), and 𝑎, 𝑏 ∇ 𝑅.

LIE GROUPS: Lie groups, named after Sophus Lie, are differentiable manifolds that carry
also the structure of a group which is such that the group operations are defined by
smooth maps. A Euclidean vector space with the group operation of vector addition is
an example of a non-compact Lie group. A simple example of a compact Lie group is the
circle: the group operation is simply rotation. This group, known as 𝑈(1), can be also
characterised as the group of complex numbers of modulus 1 with multiplication as the
group operation. Other examples of Lie groups include special groups of matrices, which
are all subgroups of the general linear group, the group of 𝑛 × 𝑛 matrices with non-zero
determinant. If the matrix entries are real numbers, this will be an 𝑛2 -dimensional
disconnected manifold. The orthogonal groups, the symmetry groups of
the sphere and hyperspheres, are 𝑛(𝑛 − 1)/2 dimensional manifolds, where 𝑛 − 1 is
the dimension of the sphere.

L’HOSPITAL’S RULE: L’ Hospital’s rule tells how to find the limit of the ratio of two
functions in cases where that ratio approaches 0/0 or 𝜘/𝜘. Let

𝑓 𝑥
𝑦=
𝑔 𝑥

Then 1’Hospital’s rule states that

𝑙𝑖𝑚 𝑥→𝑎 𝑓′ 𝑥
lim 𝑦 =
𝑙𝑖𝑚 𝑥→𝑎 𝑔′ 𝑥

Where 𝑓′ 𝑥 𝑎𝑛𝑑 𝑔′ 𝑥 represent the derivatives of these functions with respect to 𝑥.

LIKE AND UNLIKE VECTORS: Vectors having the same direction are called like vectors
and those having opposite directions are called unlike vectors.

LIKE TERMS: Two terms said to be like terms if all parts of both terms except for the
numerical coefficients are the same. For example. The terms 23𝑎2 𝑏 3 𝑐 4 and 5𝑎2 𝑏 3 𝑐 4 are
like terms. If two like terms are added , they can be combined into one tem. For example
, the sum of two terms above is 28𝑎2 𝑏 3 𝑐 4 .

LIMIT: The limit of a function is the value that the dependent variable approaches as the
independent variable approaches some fixed value. The expression “The limit of 𝑓 𝑥 as
𝑥 approaches 𝑎" is written as

𝑙𝑖𝑚 𝑥→𝑎 𝑓 𝑥

For example:

𝑥−1 𝑥+3
𝑓 𝑥 =
𝑥−1

Is undefined if 𝑥 = 1 . However the closer that 𝑥 comes to , the closer 𝑓 𝑥 approaches

4.The formal definition of limit is : The limit of 𝑓 𝑥 𝑎𝑠 𝑥 approaches 𝑎 exists and is
equal to 1, if , for any positive number 𝜀 how even small, there exists a positive number
𝛿 such that, if 0 < 𝑥 − 𝑎 < 𝛿. Then 𝑓 𝑥 − 1 < 𝜀.

LIMIT POINT OF A FILTER: Let (𝐸, 𝑇) be a topological space and 𝐹 a filter on 𝐸. We say
that a point 𝑥 ∇ 𝐸 is a limit or limit point of 𝐹 if 𝐹 is finer than the filter 𝐵(𝑥) (the basis
of open neighborhoods), i.e. if every 𝑋 ∇ 𝐵(𝑥) contains an 𝐴 ∇ 𝐹. We then say that 𝐹
converges to 𝑥 or has 𝑥 as a limit or that 𝐹 converges or is convergent.
 A filter does not necessary converge. For example the natural filter on 𝑁 with the
discrete topology.
 A filter may have more than one limit point.
 If 𝑥 is a limit point of 𝐹, 𝑥 may belong to some 𝐴 ∇ 𝐹 or may not belong to any
𝐴 ∇ 𝐹.
 For example the filter 𝐵(𝑥) consisting of the open neighborhoods of 𝑥 converges
to 𝑥, and 𝑥 belongs to every 𝑋 ∇ 𝐵(𝑥). On the other hand, the set of open
intervals of 𝑅 all having the same left-hand end point 𝑥 is a filter which
converges to 𝑥, but 𝑥 does not belong to any member of the filter.
LINDELÖF'S THEOREM: Let Ω be a half-strip in the complex plane:

Suppose that ƒ is holomorphic (i.e. analytic on Ω and that there are

constants 𝑀, 𝐴 and 𝐵 such that

and

Then 𝑓 is bounded by 𝑀 on all of Ω:

LINDEMAN THEOREM: The number 𝜋 is transcendental.

LINEAR CODE: A linear code with length 𝑛 over 𝐹𝑞 is a vector subspace of 𝐹𝑞𝑛 .
The repetition code 𝐶 = {(𝑥, . . . , 𝑥) | 𝑥 ∇ 𝐹𝑞 } is a linear code.
LINEAR COMBINATION: A vector , r, is said to be a linear combination of the vectors a, b,
c, ……etc. if there exist scalars x, y, z…etc ,such that r= xa+ yb+ zc+……. Vectors of the
form 𝛼1 𝑣1 + 𝛼2 𝑣2 + ⋯ ⋯ + 𝛼𝑛 𝑣𝑛 for 𝛼1 , 𝛼2 , ⋯ ⋯ , 𝛼𝑛 𝜖 𝐾 are called linear combination of
𝑣1 , 𝑣2 , ⋯ ⋯ , 𝑣𝑛 .

Example: Let V=ℝ2 , 𝑣1 = 1,3 , 𝑣2 = 2,5

Then 𝛼1 𝑣1 + 𝛼2 𝑣2 = 𝛼1 + 2𝛼2 , 3𝛼1 + 5𝛼2 = (0,0) iff 𝛼1 + 2𝛼2 = 0 and 3𝛼1 + 5𝛼2 = 0.
Thus we have a pair of simultaneous equations in 𝛼1 𝑎𝑛𝑑 𝛼2 and the only solution is
𝛼1 = 𝛼2 = 0, so 𝑣1 , 𝑣2 are linearly independent. Let V=ℚ2
, 𝑣1 = 1,3 , 𝑣2 = 2,6 here the equations are 𝛼1 + 2𝛼2 = 0

And 3𝛼1 + 6𝛼2 = 0 , and there are non-zero solutions such as 𝛼1 = −2 , 𝛼2 = 1 and

So 𝑣1 , 𝑣2 are linearly dependent.

Note: 𝑣1 , 𝑣2 , ⋯ ⋯ , 𝑣𝑛 are said to be linearly dependent if and only if either 𝑣1 = 0 or for

some r , 𝑣𝑟 is a linear combination of 𝑣1 , 𝑣2 , ⋯ ⋯ , 𝑣𝑟−1 .
LINEAR DEPENDENCE AND INDEPENDENCE: Let V be a vector space over the field K.
The vectors 𝑣1 , 𝑣2 , ⋯ ⋯ , 𝑣𝑛 are said to be linearly dependent if there exist scalars
𝛼1 , 𝛼2 , ⋯ ⋯ , 𝛼𝑛 𝜖 𝐾 not all zero, such that

LINEARIZATION LEMMA: Let x be a local solution of the problem. Minimize 𝜃 x subject

to x ∇ X 0 , g(x) ≤ 0. Where X 0 is an open set in Rn , and 𝜃: X 0 → R, g: X 0 → Rm are
differentiable at x.

Let

V = i g i x = 0, and g i is concave at x ,

and let

W = i g i x = 0, and g i is concave at x ,

then the system

∆θ x) + z < 0, 𝛻g 𝑤 (x z < 0,

∆g 𝑣 (x)z ≤ 0

has no solution in Rn .

LINE OF BEST FIT: The line of best fit minimizes the sum of the squares of the deviations
between each point and the line.

LINE SEGMENT: A line segment is like a piece of a line. It consists of two endpoints and
all of the points on the straight line between those two points.

LINEAR EQUATION: A linear equation with unknown 𝑥 is an equation that can be

written in the form 𝑎𝑥 + 𝑏 = 0.For example 3𝑥 − 40 = 2 can be written as 3𝑥 − 38 = 0.
38
so this is a linear equation with the solution 𝑥 = .
3

LINEAR FACTOR: A linear factor is a factor that includes only the first power of an
unknown. For example, in the expression 𝑦 = 𝑥 − 3 𝑥 2 + 3𝑥 + 4 , 𝑡𝑕𝑒 𝑓𝑎𝑐𝑡𝑜𝑟 𝑥 −
3 is a linear factor, but the factor 𝑥 2 + 3𝑥 + 4 is a quadratic factor.
LINEAR FUNCTIONAL: A linear functional on a vector space 𝑉 is a linear mapping
𝛼: 𝑉 → ℂ (or 𝛼: 𝑉 → ℝ in the real case), i.e. 𝛼 𝑎𝑥 + 𝑏𝑦 = 𝑎𝛼 𝑥 + 𝑏𝛼 𝑦 for all
𝑥, 𝑦 ∇ 𝑉 and 𝑎, 𝑏 ∇ ℂ. In simple words, A functional Φ is said to be linear if

(i) Its domain 𝐷Φ is a linear manifold, and

(ii) Φ 𝑎𝑓 + 𝑏𝑔 = 𝑎Φ 𝑓 + 𝑏Φ 𝑔 , for all 𝑓, 𝑔, ∇ 𝐷Φ

And for any scalars 𝑎 and 𝑏.

Examples:

 Let 𝑉 = 𝐶 𝑛 and 𝑐𝑘 , 𝑘 = 1, … , 𝑛 be complex numbers. Then 𝛼((𝑥1 , … , 𝑥𝑛 )) =

𝑐1 𝑥1 + ⋯ + 𝑐𝑛 𝑥𝑛 is a linear functional.
1
 On 𝐶[0,1] a functional is given by 𝛼(𝑓) = ∫0 𝑓(𝑡) 𝑑𝑡.
 On a Hilbert space 𝐻 for any 𝑥 ∇ 𝐻 a functional 𝛼𝑥 is given by 𝛼𝑥 (𝑦) = ⟨ 𝑦, 𝑥 ⟩.

LINEAR PROGRAMMING: A linear programming problem is a problem for which you

need to choose the optimal sat of values for some variables subject to some constraints.
The goal is to maximize or minimize a function called the objective function. In a linear
programming problem, the objective function and the constraints must all be linear
functions: that is, they cannot involve variables raised to any power (other than 1), and
they cannot involve two variables being multiplied together.

Some examples of problems to which linear programming can be applied include

finding the least –cost method for producing a given product , or finding the revenue-
maximizing product mix for a production facility with several capacity limitations.

Here is an example of a linear programming problem:

Maximize 6𝑥 + 8𝑦 subject:

𝑦 ≤ 10

𝑥 + 𝑦 ≤ 15

2𝑥 + 𝑦 ≤ 25

𝑥≥0

𝑦≥0
This problem has two choice variables: 𝑥 and 𝑦 . The objective function is 6𝑥 + 8𝑦, and
there are three constraints (not counting the two no negativity constraints 𝑥 ≥ 0 and
𝑦 ≥ 0. It is customary to rewrite the
constraints so that thy contain equals sign instead of inequality signs. In order to do this
some new variables called slack variables are added. One slack variable is added for
each constraint. Here is how the problem given above looks when three slack variables
𝑠1 , 𝑠2 𝑎𝑛𝑑 𝑠3 are included.

Maximize 6𝑥 + 8𝑦subject to:

𝑦 + 𝑠1 = 10

𝑥 + 𝑦 + 𝑠2 = 15

2𝑥 + 𝑦 + 𝑠3 = 25

𝑥 ≥ 0, 𝑦 ≥ 0, 𝑠1 ≥ 0, 𝑠2 ≥ 0, 𝑠3 ≥ 0

Each slack variable represents the excess capacity associated with the corresponding
constraint.

The feasible region consists of all points that satisfy the constraints . A theorem of linear
programming stated that the optimal solution will lie at one of the corner points of the
feasible region. In this case the optimal solution is at the point 𝑥 = 5, 𝑦 = 10.

A linear programming problem with two choice variables can be solved by drawing a
graph of the feasible region, as was done above. I there are more than two variables,
however it is not possible to draw a graph, and he problem must then be solved by an
algebraic procedure, such as the simplex method.

LINEARLY INDEPENDENT AND DEPENDENT SYSTEMS OF VECTORS: A system of vectors

𝑎1 , 𝑎2,……… 𝑎𝑛 is said to be linearly dependent if there exists a system of sealers 𝑥1 , 𝑥2
,…….𝑥𝑛 not all zero such that

𝑥1 𝑎1 + 𝑥2 𝑎2 + ⋯ + 𝑥𝑛 𝑎𝑛 = 0 .

A system of vectors which is no linearly dependent is said to be linearly independent.

This 𝜶 set of vectors 𝑎1 , 𝑎2 , … , 𝑎𝑛 is said to be linearly independent if every relation of
the type 𝑥1 𝑎1 + 𝑥2 𝑎2 + ⋯ + 𝑥𝑛 𝑎𝑛 = 0 𝑖𝑚𝑝𝑙𝑖𝑒𝑠 𝑡𝑕𝑎𝑡 𝑥1 = 0, 𝑥2 = 0, … . . 𝑥𝑛 = 0.
LINEARLY INDEPENDENT VECTORS: A set of vectors a, b and c is linearly independent if
it is impossible to find three scalars 𝑚, 𝑛, and 𝑝 (not all zero) such that 𝑚𝑎 + 𝑛𝑏 + 𝑝𝑐 =
0. Two vectors clearly are not linearly independent if they are multiples of each other.

LINEARLY ORDERED: An ordering ≤ (or <) of 𝐴 is called linear or total if any two
elements of 𝐴 are comparable. The pair (𝐴, ≤) is then called a linearly ordered set.
LINEAR MANIFOLD (OR SUBSPACE): A subset of 𝑀 of a linear space is called a linear
manifold (or subspace) if

 It contains all the sums of all the vectors contained in it, i.e. if 𝑓 𝑎𝑛𝑑 𝑔 are in 𝑀
then so is 𝑓 + 𝑔, and
 It contains all the scalar multiples of all its vector, i.e. if 𝑓 is in 𝑀 then so is of for
any scalar 𝑎.

Both (i) and (ii) can be combined into a single condition by saying that if
𝑓 and 𝑔 are in 𝑀, then so is 𝑎𝑓 + 𝑏𝑔, for any scalars 𝑎, 𝑏.

LINEAR OPERATOR: A linear operator T between two normed spaces 𝑋 and 𝑌 is a

mapping 𝑇: 𝑋 → 𝑌 such that 𝑇(𝜆 𝑣 + µ 𝑢) = 𝜆 𝑇(𝑣) + µ 𝑇(𝑢). The kernel of linear
operator 𝑘𝑒𝑟𝑇 and image are defined by 𝑘𝑒𝑟𝑇 = 𝑥 ∇ 𝑋: 𝑇𝑥 = 0 , Im T = {y ∇ Y: y =

Tx, for some x ∇ X}.

LINEAR PROGRAMMING: Linear programming is the most general technique that is

widely used for the optimization (maximization or minimization ) of a function to
obtain the maximum or minimum value according to certain conditions.

LINEAR PROGRAMMING PROBLEM (LPP): The LPP in general is used to optimize

(maximum/minimize) a given linear function of variables called the objective function
which is subjected to a certain set of linear equations /inequations called the constraints
or restrictions.

LINEAR PROPERTY OF FOURIER SINE AND COSINE TRANSFORMS: If 𝐹𝑠 𝑠 and 𝐺𝑠 𝑠 are

Fourier sine transforms and 𝐹𝑐 𝑠 𝑎𝑛𝑑 𝐺𝑐 𝑠 are Fourier cosine transforms of
𝑓 𝑡 𝑎𝑛𝑑 𝑔 𝑡 respectively,

Then
𝐹𝑠 𝑎 𝑓 𝑡 + 𝑏 𝑔 𝑡 = 𝑎 𝐹𝑠 𝑠 + 𝑏 𝐺𝑠 𝑠
and
𝐹𝑐 𝑎 𝑓 𝑡 + 𝑏 𝑔 𝑡 = 𝑎 𝐹𝑐 𝑠 + 𝑏 𝐺𝑐 𝑠

LINEAR PROPERTY OF FOURIER TRANSFORMS: If 𝐹 𝑠 𝑎𝑛𝑑 𝐺 𝑠 are Fourier transforms

of 𝑓 𝑡 𝑎𝑛𝑑 𝑔 𝑡 respectively,

Then
𝐹 𝑎𝑓 𝑡 + 𝑏𝑔 𝑡 =𝑎𝐹 𝑠 +𝑏𝐺 𝑠

where 𝑎 , 𝑏 are constants

LINEAR SPACE: 𝑆 is termed a linear space (or vector space) If:

1. An operation called addition and denoted by + is defined is 𝑆, which assigns to every

two vectors f, g of 𝑆, a vecot f+g of 𝑆, called the sum of 𝑓 and g, satisfying the following
properties : for any vector 𝑓, 𝑔, 𝑕 of 𝑆,

𝑓 + 𝑔 = 𝑔 + 𝑓, (commutativity)

𝑓+𝑔 + 𝑕 =𝑓+ 𝑔+𝑕 , (associativity)

The equation 𝑔 + 𝑥 = 𝑓 has at least one solution 𝑥 in 𝑆. (possibility of subtraction)

2. A certain scalar system having been prescribed which is either the complex number
system, or the real number system (the elements of the scalar defined, i.e. to every
scalar a and every vector 𝑓, a vector of called the scalar multiple of 𝑓 by a (occasionally
also denoted 𝑓𝑎) is assigned, satisfying the properties: for any scalars 𝑎, 𝑏, 𝑐 and any
vector 𝑓, 𝑔, 𝑕,

𝑎 𝑏𝑓 = 𝑎𝑏 𝑓, (Associativity w.r.t scalar multiplication to scalars)

𝑎 𝑓 + 𝑔 = 𝑎𝑓 + 𝑎𝑔, (Distributivity of multiplication over addition)

1𝑓 = 𝑓 (property of multiplication by unity)

The linear space is called a complex linear space or a real linear space according as the
prescribed scalar system is the complex number system or the real number system.

LINEAR SPAN: Let 𝐴 be a subset (finite or infinite) of a normed space 𝑉. The linear span
of 𝐴, write 𝐿(𝐴), is the intersection of all linear subspaces of 𝑉 containing 𝐴, i.e. the
smallest subspace containing 𝐴, equivalently the set of all finite linear combination of
elements of 𝐴. The closed linear span of 𝐴 write 𝐶𝐿(𝐴) is the intersection of all closed
linear subspaces of 𝑉 containing 𝐴, i.e. the smallest closed subspace containing 𝐴.

 If A is a subset of finite dimension space then 𝐿(𝐴) = 𝐶𝐿(𝐴).

 For an infinite A spaces L(A) and CL(A) could be different.
 If H is an inner product space and sequences 𝑥𝑛 and 𝑦𝑛 have limits x and y
correspondingly. Then ⟨ 𝑥𝑛 , 𝑦𝑛 ⟩ → ⟨ 𝑥, 𝑦 ⟩.

LINEAR TRANSFORMATIONS: It is important to study special classes of functions as it is

to study special classes of objects. Often these are functions special classes of certain
properties or structures. For example, continuous functions preserves which points are
close to which other points. In linear algebra, the functions, which preserve the vector
space structures, are called linear transformations.

Let U and V be two vector spaces over the field K. A linear transformation or linear map
T from U to V is a function T: U→ 𝑉 𝑠𝑢𝑐𝑕 𝑡𝑕𝑎𝑡

(i) 𝑇 𝑢1 + 𝑢2 = 𝑇 𝑢1 + 𝑇 𝑢2 for all 𝑢1 , 𝑢2 𝜖 𝑈

(ii) T(𝛼𝑢) = 𝛼𝑇 𝑢 𝑓𝑜𝑟 𝑎𝑙𝑙 𝛼 𝜖 𝐾, 𝑢 𝜖 𝑈
The above two conditions for linearity are equivalent to a single condition

𝑇 𝛼𝑢1 + 𝛽𝑢2 = 𝛼𝑇 𝑢1 + 𝛽𝑇 𝑢2 for all 𝑢1 , 𝑢2 𝜖 𝑈, 𝛼, 𝛽 𝜖 𝐾

Let 𝑇: 𝑈 → 𝑉 be a linear map.Then

(i) 𝑇(0𝑣) = 0𝑣
(ii) 𝑇(−𝑢) = −𝑇(𝑢) for all u 𝜖 𝑈
Examples: Many familiar geometrical transformations, such as projections, rotations,
reflections and magnifications are linear maps. Note, however, that a nontrivial
translation is not a linear map, because it does not satisfy 𝑇(0𝑣) = 0𝑣.

1. Let 𝑈 = ℝ3 , 𝑉 = ℝ2 and define 𝑇: 𝑈 → 𝑉 by 𝑇((𝛼, 𝛽, 𝛾)) = (𝛼, 𝛽).Then T is a linear

map. This type of mapping is known as projection.

2. U = V = ℝ2 .we interpret v in ℝ2 as a directed line vector from 0 to 𝑣 and let 𝑇(𝑣) be

the vector obtained by rotating 𝑣 through an angle 𝜃 anticlockwise about the origin.It is
easy to see that 𝑇 𝑢1 + 𝑢2 = 𝑇 𝑢1 + 𝑇 𝑢2 𝑎𝑛𝑑 T(𝛼𝑢) = 𝛼𝑇 𝑢 and so T is a linear
map.By considering the unit vectors, we have 𝑇(1,0) = (𝑐𝑜𝑠 𝜃, 𝑠𝑖𝑛 𝜃) 𝑎𝑛𝑑 𝑇(0,1) =
(−𝑠𝑖𝑛 𝜃, 𝑐𝑜𝑠 𝜃) and hence

𝑇(𝛼, 𝛽) = 𝛼 𝑇(1,0) + 𝛽𝑇 0,1 = (𝛼 𝑐𝑜𝑠 𝜃 − 𝛽 𝑠𝑖𝑛 𝜃, 𝛼 𝑠𝑖𝑛 𝜃 + 𝛽 𝑐𝑜𝑠 𝜃)

3. Let us take U = V = ℝ2 . Let 𝑇(𝑣) be the vector obtained by resulting from reflecting
𝜃
𝑣 through a line through the origin that makes an angle 2 with the 𝑥 −axis. This is again

a linear map. We find that 𝑇(1,0) = (𝑐𝑜𝑠 𝜃, 𝑠𝑖𝑛 𝜃) 𝑎𝑛𝑑 𝑇(0,1) = (−𝑠𝑖𝑛 𝜃, 𝑐𝑜𝑠 𝜃) and
𝑇(𝛼, 𝛽) = 𝛼 𝑇(1,0) + 𝛽T 0,1 = (α cos θ + β sin θ, α sin θ − β cos θ)

4. Let 𝑈 = 𝑉 = ℝ 𝑥 , the set of polynomials over ℝ, and let T be differentiation;i.e

T(p(x))=𝑝′ 𝑥 𝑓𝑜𝑟 𝑝 𝑥 𝜖 ℝ 𝑥 .This is also a linear map.

5. Let U= 𝐾 𝑥 , the set of polynomials over K. Every 𝛼 𝜖 𝐾 gives rise to two linear maps,
shift 𝑆𝛼 : 𝑈 → 𝑈, 𝑆𝛼 𝑓 𝑥 = 𝑓 𝑥 − 𝛼 and evaluation 𝐸𝛼 : 𝑈 → 𝐾, 𝐸𝛼 𝑓 𝑥 =𝑓 𝛼

6. For any vector space V, we define the identity map

𝐼𝑣 : 𝑉 → 𝑉 𝑏𝑦 𝐼𝑣 𝑣 = 𝑣 𝑓𝑜𝑟 𝑎𝑙𝑙 𝑣 𝜖 𝑉. This is a linear map.

7. For any vector space U, V over the field K, we define the map 0𝑈,𝑉 : U→ 𝑉 by 0𝑈,𝑉
(u)=0v for all 𝑢 𝜖 𝑈. This is also a linear map.

Note: Let U, V be vector spaces over K. let 𝑢1 , 𝑢2 , ⋯ ⋯ , 𝑢𝑛 be a basis of U and let

𝑣1 , 𝑣2 , ⋯ ⋯ , 𝑣𝑛 be any sequence of n vectors in V.Then there is a unique linear map

T:U→ 𝑉 with T(𝑢𝑖 ) = 𝑣𝑖 for 1 ≤ 𝑖 ≤ 𝑛.

To any linear map 𝑇: 𝑈 → 𝑉,we can associate a special subspace of U and a special
subspace of V called Kernel and Range spaces.

Let 𝑇: 𝑈 → 𝑉 be a linear map.The image of 𝑇, written as 𝐼𝑚 𝑇 , is the set of vectors 𝑣 𝜖 𝑉

such that 𝑣 = 𝑇(𝑢) for some 𝑢 𝜖 𝑈, is called Range of 𝑇.

The kernel of T, written as 𝐾𝑒𝑟(𝑇), is the set of vectors 𝑢 𝜖 𝑈 such that 𝑇(𝑢) = 0𝑣

𝐼𝑚 𝑇 = 𝑇 𝑢 : 𝑢 𝜖 𝑈 ; 𝐾𝑒𝑟(𝑇) = 𝑢 𝜖 𝑈: T(u) = 0v

Examples: Let us consider the examples 1-7 above

  In example 1, 𝐾𝑒𝑟(𝑇) = {(0,0, 𝛾): 𝛾 𝜖 𝑅} and 𝐼𝑚(𝑇) = ℝ2
 In example 2 and 3, 𝐾𝑒𝑟(𝑇) = {0} and 𝐼𝑚(𝑇) = ℝ2
 In example 4, 𝐾𝑒𝑟(𝑇)is the set of constant polynomials and 𝐼𝑚(𝑇) = ℝ 𝑥
 In example 5, 𝐾𝑒𝑟 (𝑆𝛼 ) = 0 and 𝐼𝑚(𝑆𝛼 ) = 𝐾 𝑥 ,while 𝐾𝑒𝑟 𝐸𝛼 is the set of all
polynomials divisible by 𝑥 − 𝛼 , and 𝐼𝑚 𝐸𝛼 = 𝐾.
 In example 6, 𝐾𝑒𝑟 (𝐼𝑣 ) = 0 and 𝐼𝑚(𝐼𝑣 ) = 𝑉
 In example 7, 𝐾𝑒𝑟 (0𝑈,𝑉 ) = 𝑈 and 𝐼𝑚(0𝑈,𝑉 ) = 0
Let 𝑇: 𝑈 → 𝑉 be a linear map.Then

(i) 𝐼𝑚(𝑇) is a subspace of 𝑉.

(ii) 𝐾𝑒𝑟(𝑇) is a subspace of 𝑈.
The dimensions of the kernel and image of a linear map contain important information
about it, and are related to each other

Let 𝑇: 𝑈 → 𝑉 be a linear map

(i) 𝐷𝑖𝑚(𝐼𝑚(𝑇)) is called the rank of T.

(ii) 𝐷𝑖𝑚(𝐾𝑒𝑟(𝑇)) is called the nullity of T
LINEAR TRANSFORMATIONS AND MATRICES: Let 𝑇: 𝑈 → 𝑉 be a linear map, where
𝑑𝑖𝑚(𝑈) = 𝑛 , 𝑑𝑖𝑚(𝑉) = 𝑚.Suppose that we are given a basis 𝑒1 , 𝑒2 , ⋯ ⋯ , 𝑒𝑛 of 𝑈 and a
basis 𝑓1 , 𝑓2 , ⋯ ⋯ , 𝑓𝑚 of 𝑉. Now, for 1 ≤ 𝑗 ≤ 𝑛 , the vector 𝑇(𝑒𝑗 ) lies in 𝑉,so 𝑇(𝑒𝑗 )can be
written uniquely as a linear combination of 𝑓1 , 𝑓2 , ⋯ ⋯ , 𝑓𝑚 . Let

𝑇 𝑒1 = 𝛼11 𝑓1 + 𝛼21 𝑓2 + ⋯ ⋯ ⋯ + 𝛼𝑚1 𝑓𝑚

𝑇 𝑒2 = 𝛼12 𝑓1 + 𝛼22 𝑓2 + ⋯ ⋯ ⋯ + 𝛼𝑚2 𝑓𝑚

⋯⋯⋯⋯⋯⋯⋯⋯⋯⋯⋯⋯⋯⋯⋯⋯⋯

⋯⋯⋯⋯⋯⋯⋯⋯⋯⋯⋯⋯⋯⋯⋯⋯⋯

𝑇 𝑒𝑛 = 𝛼1𝑛 𝑓1 + 𝛼2𝑛 𝑓2 + ⋯ ⋯ ⋯ + 𝛼𝑚𝑛 𝑓𝑚

Where the coefficients 𝛼𝑖𝑗 𝜖 𝐾 ( 1 ≤ 𝑖 ≤ 𝑚 ,1 ≤ 𝑗 ≤ 𝑛) are uniquely determined.

Putting it more compactly ,we define scalars 𝛼𝑖𝑗 by

𝑚
𝑇 𝑒𝑗 = 𝑖=1 𝛼𝑖𝑗 𝑓𝑖 for 1 ≤ 𝑗 ≤ 𝑛.
The coefficients 𝛼𝑖𝑗 form an 𝑚 × 𝑛 matrix

𝛼11 𝛼12 ⋯ 𝛼1𝑛

𝛼21 𝛼22 ⋯ 𝛼2𝑛
𝐴= ⋮ ⋮ ⋮ ⋮
𝛼𝑚1 𝛼𝑚2 ⋯ 𝛼 𝑚𝑛

over K. Then A is called the matrix of the linear map T with respect to the chosen bases
of U and V. In general, different choices of bases give different matrices.

Note the role of individual columns in A. The jth column of A consists of the coordinates
of 𝑇 𝑒𝑗 w.r.t the basis 𝑓1 , 𝑓2 , ⋯ ⋯ , 𝑓𝑚 of V.

LINE OF CURVATURE: A curve on a surface is called a line of curvature if the tangent at

any point of it is along the principal direction at that point.

LINE OF FLOW: A line of flow is a line whose direction coincides with the direction of
the resultant velocity of the fluid.

LINGUISTIC VARIABLE: A linguistic variable is characterized by a quintuple

(𝑥, 𝑇(𝑥), 𝑈, 𝐺, ˜𝑀), in which 𝑥 is the name of the variable, 𝑇(𝑥) (or simply 𝑇) denotes the
term set of 𝑥, that is, the set of names of linguistic values of 𝑥. Each of these values is a
fuzzy variable, denoted generically by 𝑋 and ranging over a universe of discourse 𝑈,
which is associated with the base variable 𝑢; 𝐺 is a syntactic rule (which usually has the
form of a grammar) for generating the name, 𝑋, of values of 𝑥. 𝑀 is a semantic rule for
associating with each 𝑋 its meaning. ˜𝑀(𝑋) is a fuzzy subset of 𝑈. A particular 𝑋, that is,
a name generated by 𝐺, is called a term
LIOUVILLE’S FORMULA FOR 𝜿𝒈 (DIFFERENTIAL GEOMETRY): If a given curve makes an
angle 𝜃 with the parametic curve v=c, then according to Lioville’s formula, 𝜅𝑔 is
expressed by

𝜅𝑔 = 𝜃 ′ + 𝑃𝑢′ + 𝑄𝑣′

1
Where 𝑃 = 2𝐸𝐹1 − 𝐹𝐸1 − 𝐸𝐸2
2𝐻𝐸

1
and 𝑄 = (𝐸𝐺 − 𝐹𝐸2 )
2𝐻𝐸
LIOUVILLE’S THEOREM COMPLEX ANALYSIS): If an entire function is bounded for all
values of 𝑧, then it is constant.

or, If a function 𝑓(𝑧) is analytic for finite values of 𝑧, and is bounded, then 𝑓(𝑧) is
constant.

or, if 𝑓 is regular in whole 𝑧 − 𝑝𝑙𝑎𝑛𝑒 𝑎𝑛𝑑 𝑖𝑓 𝑓(𝑧) < 𝑘∀𝑧, then 𝑓(𝑧) must be constant.

or, Every holomorphic function 𝑓 for which there exists a positive number 𝑀 such that
|𝑓(𝑧)| ≤ 𝑀 for all 𝑧 in 𝐶 is constant.

LIOUVILLE'S THEOREM (CONFORMAL MAPPING): Any smooth conformal mapping on a

domain of 𝑹𝒏 , where 𝑛 > 2, can be expressed as a composition
of translations, similarities, orthogonal transformations and inversions: they are
all Möbius transformations. This severely limits the variety of possible conformal
mappings in 𝑹𝟑 and higher-dimensional spaces. By contrast, conformal mappings
in 𝑹𝟐 can be much more complicated – for example, all simply connected planar
domains are conformally equivalent, by the Riemann mapping theorem.
LIOUVILLE'S THEOREM (HAMILTONION MECHANICS): The phase-space distribution
function is constant along the trajectories of the system — that is that the density of
system points in the vicinity of a given system point travelling through phase-space is
constant with time.
LIPSCHITZ CONDITION: A function 𝑓(𝑥, 𝑦) continuous on a rectangle 𝑅 satisfies a
Lipschitz condition with constant 𝐴 if there exists a real positive 𝐴 such that
|𝑓(𝑥, 𝑢) − 𝑓(𝑥, 𝑣)| ≤ 𝐴|𝑢 − 𝑣| 𝑓𝑜𝑟 (𝑥, 𝑢) ∇ 𝑅, (𝑥, 𝑣) ∇ 𝑅.
This is a new condition on a function, stronger than being continuous but weaker than
being differentiable.
LITERAL NUMBER: A literal number is a number expressed as a numeral, not as a
variable. For example, in the equation 𝑥 = 34𝑦, 3.4 is a literal number.

LITTLE PICARD THEOREM: If a function 𝑓 : 𝑪 → 𝑪 is entire and non-constant, then the

set of values that 𝑓(𝑧) assumes is either the whole complex plane or the plane minus a
single point.
LOBACHEVSKY: Nikolay Lobachevsky (1792 to 1856) was a Russian mathematician
who developed a version of non- Euclidian geometry.
LOCAL CONVECTIVE AND MATERIAL DERIVATIVES: Let a fluid element moves from the
point 𝑟 in a flow field which is uniform through the space at any instant. The change in
any fluid property 𝐹 for the element in a small interval 𝛿𝑡 is 𝜕𝐹/𝛿𝑡 𝜕𝑡, where 𝜕𝐹/𝛿𝑡
represents the rate at which 𝐹 is changing locally at the point 𝑟. This change is known as
the local derivative or local rate of change due to temporal changes.

Let a fluid element moves in a steady and non-uniform flow field from a point 𝑟 to a new
position (𝑟 + 𝑞𝛿𝑡). Then the change in any fluid property 𝐹 is equal to (𝜕𝐹/𝜕𝑆)𝑞𝛿𝑡 (to
the first app. in 𝑞𝛿𝑡). This change due to movement of the element is known as the
convective change or the convective derivative.

The change in any fluid property 𝐹 for a fluid element moving in the unsteady and non-
∂F
uniform flow field is made up of both the local and convective changes, i.e., 𝛿𝑡 +
∂t
𝜕𝐹
𝑞𝛿𝑡 , which is the same as the rate of change of 𝐹 following a certain fluid element.
𝜕𝑆

This is known a substantial derivative or a material derivative.

LOCAL DOMAIN: A local domain is an integral domain that has exactly one maximal
ideal. A local domain is a discrete valuation ring if and only if it is a principal ideal
domain.

LOCALISED AND FREE VECTORS: A vector which is drawn parallel to a given vector
through a specified point in space is called a localized vector. There can be one and only
one such vector. But if the origin of vector is not specified , the vectors are said to be
free vectors.

LOCALLY FINITE: A collection of subsets of a topological space 𝑋 is said to be locally

finite, if for each element 𝑥 ∇ 𝑋 there exists a neighborhood which intersects non-
trivially with only finitely many of the subsets.
LOCALLY FINITE ATLAS: An atlas of an abstract manifold 𝑀 is said to be locally finite if
the collection of images 𝜍(𝑈) is locally finite in 𝑀.
LOCALLY ISOMETRIC: Two surface S and S’ are called locally isometric if every point of
the surface S has a neighbourhood which is isometric with a region of S’.
LOCALLY PATHWISE CONNECTED: A topological space 𝑋 is said to be locally pathwise
connected if it has the following property. For each point 𝑥 ∇ 𝑋 and each neighborhood
𝑉 there exists an open pathwise connected set 𝑈 such that
𝑥 ∇ 𝑈 ⊂ 𝑉.
LOCAL MAXIMUM: A local maximum point for a function 𝑦 = 𝑓 𝑥 is a point where the
value of 𝑦 is larger than the points near it. If the first derivative is zero and the second
derivative is negative at a point 𝑥1 , 𝑓 𝑥1 , then the function has a local maximum at
that point . There may be more than one local maximum, so there is no guarantee that a
particular local maximum will be the absolute maximum.

LOCAL MINIMUM: A local minimum point for a function 𝑦 = 𝑓 𝑥 is a point where the
value of y is smaller than the points near it. If the first derivative is zero and the second
derivative is negative at a point 𝑥1 , 𝑓 𝑥1 , then the function has a local maximum at
that point . There may be more than one local minimum, so there is no guarantee that a
particular local minimum will be the absolute minimum.

LOCUS: The term “locus” is a technical way of saying “set of points”. For example, a circle
can be defined as being “the locus of points in a plane that are a fixed distance from a
given point”. The plural of “locus” is “loci”

LOG: The function y= log x is an abbreviation for the logarithm function to the base 10.

LOGARITHM: The equation 𝑥 = 𝑎 𝑦 can be written as 𝑦 = 𝑙𝑜𝑔𝑎 𝑥, which means "𝑦 is the
logarithm to the base 𝑎 of 𝑥 and 𝑦 is the exponent to which 𝑎 must be raised in order to
result in 𝑥. For example 𝑙𝑜𝑔2 8 = 3 means the same as 23 = 8". Any positive number
(except1) can be used as the base for a logarithm function. The two most useful bases
are 10 and 𝑒. Logarithms to the base 10 are called common logarithms.

Logarithm form exponential form

𝑙𝑜𝑔1 = 0 100 = 1

𝑙𝑜𝑔10 = 1 101 = 10

𝑙𝑜𝑔100 = 2 102 = 100

𝑙𝑜𝑔1,000 = 3 103 = 1,000

Except in a few simple cases, Logarithms will be irrational numbers.

Logarithms satisfy these properties:

𝑙𝑜𝑔 𝑥𝑦 = log 𝑥 + log 𝑦

𝑙𝑜𝑔 𝑦/𝑥 = 𝑙𝑜𝑔𝑦 − log 𝑥

𝑙𝑜𝑔 𝑥 𝑛 = 𝑛 𝑙𝑜𝑔 𝑥

Logarithms have also been very helpful as calculation aids, because a multiplication
problem can be turned into an addition problem by taking the logarithms.

LOGARITHMIC TEST: The series 𝑢𝑛 of positive terms is convergent or divergent

𝑢𝑛
according as lim𝑛→∞ 𝑛 𝑙𝑜𝑔 > 1 𝑜𝑟 < 1.
𝑢 𝑛 +1

LOGIC: Logic is the study of sound reasoning. The study of logic focuses on the study of
arguments. An argument is a sequence of sentences (called premises), that lead to a
resulting sentences (called the conclusion). An argument is a valid argument if the
conclusion does follow from the premises. In order words, if an argument is valid and all
its premises are true, then the conclusion must be true.

Logic can be used to determine whether an argument is valid: however, logic alone
cannot determine whether the premises are true or false. Once an argument has been
shown to be valid, then all other arguments of the same general form will also be valid,
even if their premises are different.

Arguments are composed of sentences. Sentences are said to have the truth value 𝑇
corresponding to what we normally think of as “true” or the truth value
𝐹 corresponding to “false” . In studying the general logical properties of sentences, it is
customary to represent a sentence by a lower-case letter, such as p, q or r called a
sentence variable or a Boolean variable. Sentences either can be simple sentences or can
consist of simple sentences joined by connectives and called compound sentences.

LOGICAL SYMBOLS: If A and B are propositions, the propositions (A and B), (A or B), (A
implies B), and (not A) are denoted by 𝐴 ⋀𝐵, 𝐴⋁𝐵, 𝐴 ⟶ 𝐵, ∼ 𝐴, respectively. We call ∼
the negation of ∼ 𝐴, 𝐴 ⋀𝐵 the conjunction (or logical product), 𝐴⋁𝐵 the disjunction (or
logical sum), and 𝐴 ⟶ 𝐵 the implication (or B by A). The proposition (𝐴 ⟶ 𝐵)⋀(𝐵 ⟶
𝐴) is denoted by 𝐴 ⟷ 𝐵 and is read “A and B are equivalent.” 𝐴⋁ 𝐵 means that at least
one of A and B holds.
LORENTZ TRANSFORMATION: The Lorentz transformation describes how events look
different to observes moving with different velocities, according to Einstein’s special
theory of relativity . let t, x, y, z be the time and space coordinates of an event in the
original frame of reference , and let t’ ,x’ ,y’ ,z’ be the coordinates of an event in a new
frame, which is moving with velocity v in the positive x direction with respect to the
original frame.

1
Define 𝛾 = 1−𝑣 2/ 𝑐 2

Where 𝑐 is the speed of light.

Then the Lorentz transformation can be written:

Time coordinates Space coordinates

𝑡 𝑥 𝑦 𝑧 𝑜𝑟𝑖𝑔𝑖𝑛𝑎𝑙 𝑓𝑟𝑎𝑚𝑒
𝑣𝑥 𝑥 ′ = 𝛾 𝑥 − 𝑣𝑡 𝑕′ = 𝑦 𝑧′ = 𝑧 New frame
𝑡′ = 𝛾 1 +
𝑐2

In everyday life, the velocity 𝛾is always very small compared to the speed of light. So 𝛾 is
always very close to 1.

LOXODROME: A loxodrome on a sphere is curve that makes a constant angle with the
parallels of latitude.

𝛼1 𝑣1 + 𝛼2 𝑣2 + ⋯ ⋯ + 𝛼𝑛 𝑣𝑛 = 0

And 𝑣1 , 𝑣2 , ⋯ ⋯ , 𝑣𝑛 are said to be linearly independent if they are not linearly

dependent.In other words they are linearly independent if the only scalars
𝛼1 , 𝛼2 , ⋯ ⋯ , 𝛼𝑛 𝜖 𝐾 that satisfy the above equations are

𝛼1 = 0, 𝛼2 = 0, ⋯ ⋯ , 𝛼𝑛 = 0.
LOWER RIEMANN INTEGRAL: Let 𝑓 be a bounded function on a bounded interval [𝑎, 𝑏].
Then lower Riemann integral of f over [𝑎, 𝑏] is the supremum of 𝐿 𝑃, 𝑓 over all
𝑏
partitions 𝑃 of [𝑎, 𝑏] and is denoted as ∫𝑎 𝑓𝑑𝑥.

LOWER RIEMANN-STIELTJES SUM: Let 𝑓 be a bounded real valued function on a

bounded interval [𝑎, 𝑏] and let 𝑔 be a monotonically non-decreasing function on [𝑎, 𝑏].
Let 𝑃 = 𝑥0 , 𝑥1 , 𝑥2 , 𝑥3 , … , 𝑥𝑛 be a partition of [𝑎, 𝑏]. We write
∆𝑔𝑖 = 𝑔𝑖 − 𝑔𝑖−1
Then lower Riemann-Stieltjes sum denoted by 𝐿(𝑃, 𝑓, 𝑔) is defined as
𝑛

𝐿 𝑃, 𝑓, 𝑔 = 𝑚𝑖 ∆𝑔𝑖
𝑖=1

Where
𝑚𝑖 = sup 𝑓(𝑥)
𝑥∇[𝑥 𝑖−1 ,𝑥 𝑖 ]

LOWER RIEMANN SUM: Let 𝑓 be a bounded function on a bounded interval [𝑎, 𝑏]. Let
𝑃 = 𝑥0 , 𝑥1 , 𝑥2 , 𝑥3 , … , 𝑥𝑛 be a partition of [𝑎, 𝑏]. Then lower Riemann sum denoted by
𝐿(𝑃, 𝑓) is defined as
𝑛

𝐿 𝑃, 𝑓 = 𝑚𝑖 𝑥𝑖 − 𝑥𝑖−1
𝑖=1

Where
𝑚𝑖 = sup 𝑓(𝑥)
𝑥∇[𝑥 𝑖−1 ,𝑥 𝑖 ]

LOWER TRIANGULAR MATRIX: A square matrix A = aij is called an lower triangular

matrix if aij = 0 whenever i < 𝑗.

Thus in an lower triangular matrix all the elements above the principle diagonal are
zero.

a11 0 0 …0
a21 a22 0 …0
For example A = a31 a32 a33 …0
… … … ……
an1 an2 an3 … ann n×n

is an lower triangular matrix of the size n × n. Similarly

 3 0 0 0
1 0 4 −1 0 0
A= ,B =
2 3 2×2 0 2 0 0
5 7 1 2 4×4

are the lower triangular matrices.

A triangular matrix A = aij is called strictly triangular if aij = 0 for all i = 1,2, … , n.
n×n

LOWSEST COST ENTRY (MINIMA) MEHOD: This method of solution can be taken similar
to the North-West corner rule. But the difference lies in the fact that in this case, we first
assign allocations to the cell having lowest value. But in case , if either two or three cells
have the same least value, then the allocation is allotted to that cell for which the value
of demand is also least. In this way, we proceed from most minimum to the next
minimum.

LOWEST UPPER BOUND: Let 𝑆 be a set with an ordering relation ≤, and let 𝑇 be a subset
of 𝑆. A lowest upper bound of 𝑇 is an upper bound 𝑥 of 𝑇 with the property that 𝑥 ≤ 𝑦
for every upper bound 𝑦 of 𝑇. A lowest upper bound of 𝑇, when it exists, is unique.
LUZIN’S THEOREM: Every bounded measurable functions on [a,b] is the limit in the
mean of p-th power of a sequence of continuous function i.e., if 𝑓 is a bounded
measurable function on 𝑎, 𝑏 , then for a given 𝜀 > 0, there exists a continuous function
𝑔 on [𝑎, 𝑏], such that

𝑓−𝑔 𝑝 < 𝜀.

MACLAURIN: Colin Maclaurin (1698 to1746) was a Scottish mathematician who

extended the field of calculus.

MACLAURIN SERIES: The Maclaurin series is a special case of the Taylor series for
𝑓 𝑥 + 𝑕 , 𝑤𝑕𝑒𝑛 𝑥 = 0.

MACLAURIN’S THEOREM WITH CAUCHY’S FORM OF REMAINDER: If 𝑓(𝑥) is a single

valued function of 𝑥 such that
  all the derivatives of 𝑓(𝑥) upto (𝑛 − 1)th order are continuous in the closed
interval [0, 𝑥]
 𝑓 (𝑛 ) (𝑥) exists in (0, 𝑥)
then
𝑥 2 ′′ 𝑥 𝑛 −1
𝑓 𝑥 = 𝑓 0 + 𝑥𝑓’ 0 + 𝑓 0 +−−−+ 𝑓 𝑛−1 0
2! (𝑛 − 1)!
𝑥𝑛 𝑛−1 𝑛
+ 1−𝜃 𝑓 𝜃𝑥
(𝑛 − 1)!
where 𝜃 ∇ 0, 1
MACLAURIN’S THEOREM WITH LAGRANGE’S FORM OF REMAINDER: If 𝑓(𝑥) is a single
valued function of 𝑥 such that
 all the derivatives of 𝑓(𝑥) upto (𝑛 − 1)th order are continuous in the closed
interval [0, 𝑥]
 𝑓 (𝑛 ) (𝑥) exists in (0, 𝑥)
then
𝑥 2 ′′ 𝑥 𝑛−1 𝑥𝑛
𝑓 𝑥 = 𝑓 0 + 𝑥𝑓’ 0 + 𝑓 0 +−−−+ 𝑓 𝑛−1 0 + 𝑓 𝑛 𝜃𝑥
2! (𝑛 − 1)! 𝑛!
where 𝜃 ∇ 0, 1
MAGNIFICATION OR STRETCHING (COMPLEX ANALYSIS): Consider the map 𝑤 = 𝑎𝑧,
where 𝛼 is real. The two figures in 𝑧 − 𝑝𝑙𝑎𝑛𝑒 and 𝑤 − 𝑝𝑙𝑎𝑛𝑒 are similar and similarly
situated about their respective origins, but the figure in 𝑤 − 𝑝𝑙𝑎𝑛𝑒 is a times the figure
in 𝑧 − 𝑝𝑙𝑎𝑛𝑒. Such map is called magnification or stretching or dilation.

MAGNITUDE OF A VECTOR: The magnitude of a vector a is its length. It is symbolized by

two pairs of vertical lines, and it can be found by taking the square root of the dot
product of he vector with itself:

𝑎 = 𝑎∙𝑎

MAINARDI-CODAZZI EQUATIONS (DIFFERENTIAL GEOMETRY): The equations:

𝐿2 − 𝑀1 = 𝑚𝐿 − 𝑙 − 𝜇 𝑀 − 𝜆𝑁

and

𝑀2 − 𝑁1 = 𝑛𝐿 − 𝑚 − 𝑣 𝑀 − 𝜇𝑁
are called Mainardi-Codazzi equations.

MAJOR AXIS: The major axis of an ellipse is the line segment joining two points on the
ellipse that passes through the two foci. It is the longest possible distance across the
ellipse.

MANDELBROT SET: The Mandelbrot set , discovered by Benoit Mandelbrot , is a famous

fractal i.e., a shape containing an infinite amount of fine detail . It is the set of values of c
for which the series 𝑧𝑛+1 = 𝑧𝑛2 + 𝑐 converges, where 𝑧 𝑎𝑛𝑑 𝑐 are complex numbers and
𝑧 is initially (0,0).

MANIFOLD: A manifold is a topological space that resembles euclidean spacenear each

point. More precisely, each point of an n-dimensional manifold has aneighbourhood that
is homeomorphic to the euclidean space of dimension n. Linesand circles, but not figure
eights, are one-dimensional manifolds. Two-dimensional manifolds are also
called surfaces. Examples include the plane, the sphere, and thetorus, which can all be
realized in three dimensions, but also the klein bottle andreal projective plane which
cannot.
A topological manifold is a locally Euclidean Hausdorff space. A manifold need not
be paracompact or second-countable. A Ck manifold is a differentiable manifold whose
chart overlap functions are k times continuously differentiable. A C∞ or smooth
manifold is a differentiable manifold whose chart overlap functions are infinitely
continuously differentiable. Some spcial types of manifolds are

 A complex manifold is a manifold modeled on ℂ𝑛 with holomorphic transition

functions on chart overlaps. These manifolds are the basic objects of study in
complex geometry. A one-complex-dimensional manifold is called a Riemann
surface. Note that an n-dimensional complex manifold has dimension 2n as a real
differentiable manifold.
 A 'CR manifold' is a manifold modeled on boundaries of domains in .
 Banach manifolds which are locally homeomorphic to Banach spaces. Similarly,
Fréchet manifolds are locally homeomorphic to Fréchet spaces.
 A symplectic manifold is a kind of manifold which is used to represent the phase
spaces in classical mechanics. They are endowed with a 2-form that defines
the Poisson bracket. A closely related type of manifold is a contact manifold.
  A 'combinatorial manifold' is a kind of manifold which is discretization of a
manifold. It usually means a piecewise linear manifold made by simplicial
complexes.
 A 'digital manifold' is a special kind of combinatorial manifold which is defined in
digital space. See digital topology.

MANIFOLD WITH BOUNDARY: A manifold with boundary is a manifold with an edge, e.g.
a sheet of paper is a 2-manifold with a 1-dimensional boundary. The boundary of an 𝑛-
manifold with boundary is an (𝑛 − 1) manifold. A disk is a 2-manifold with boundary.
Its boundary is a circle, a 1-manifold. A square with interior is also a 2-manifold with
boundary. A closed ball is a 3-manifold with boundary. Its boundary is a sphere, a 2-
manifold.

In technical language, a manifold with boundary is a space containing both interior

points and boundary points. Every interior point has a neighborhood homeomorphic to
the open n-ball {(𝑥1 , 𝑥2 , … , 𝑥𝑛 ) | 𝑥𝑖2 < 1}. Every boundary point has a neighborhood
homeomorphic to the "half" n-ball {(𝑥1 , 𝑥2 , … , 𝑥𝑛 ) | 𝑥𝑖2 < 1 𝑎𝑛𝑑 𝑥1 ≥ 0}. The
homeomorphism must send each boundary point to a point with 𝑥1 = 0.

MANTISSA: The mantissa is the part of a common logarithm to the right of the decimal
point. For example, in the expression log 115=2.0607, the quantity 0.0607 is the
mantissa.

MAPPING: A mapping is a rule that, to each member of one set, assigns a unique member
of another set.

MAPPINGS OF SURFACE (DIFFERENTIAL GEOMETRY): The two surface are said to be

equivalent, if this correspondence of points between two surfaces preserves some
geometrical structure on the surface and these are said to be mapped upon one another
if there exists a one-to-one correspondence between the points of the two surface. Let S
and S*, the two surfaces, have the parametric representations

𝑟 = 𝑟 𝑢, 𝑣 , 𝑟 ∗ = 𝑟 ∗ (𝑢∗ , 𝑣 ∗ )

respectively. The most general map of S and S* is defined by equations.

𝑢∗ = 𝜙 𝑢, 𝑣 , 𝑣 ∗ = Ψ (𝑢, 𝑣)
This establishes a one-to-one correspondence between the permissible pair of values of
𝑢, 𝑣 and the permissible pairs of values of 𝑢∗ , 𝑣 ∗ where if is assumed that 𝜙 𝑎𝑛𝑑 Ψ are
analytic functions of u and v. The parameters u,v on the surface S may be introduced as
parameters on the surface 𝑆 ∗ by putting the values of 𝑢∗ , 𝑣 ∗ , in terms of 𝑢, 𝑣 in
𝑟 ∗ = 𝑟 ∗ 𝑢∗ , 𝑣 ∗ . Thus the corresponding points of the two surface have the same
curvilinear co-ordinates. Thus the parametric equations of the two surface can be
written as

𝑟 = 𝑟 𝑢, 𝑣 𝑎𝑛𝑑 𝑟 ∗ (𝑢, 𝑣)

These equations serve not only to define the surface but also the map established
between them.

MATERIAL LINE ELEMENT: A material line element is a small line element marked in
the fluid, i.e., made up of fluid parcels, with end points moving with the flow.

MATHEMATICAL INDUCTION: Mathematical induction is a method for proving that a

proposition is true for 𝑛 = 1 . Then show that, if the proposition is true for an arbitrary
number k then it must be true for next number 𝑘 + 1. Once we have done these two
steps, the proposition has been proved, since, if it is true for 1, then it must also be true
for 2, which means it must be true for 3, which means it must be true for 4, and so on.

MATRIX: A set of mn numbers (real or complex) arranged in the form of a rectangular

array having m row having m rows and n columns is called an m × n matrix [to be read
as m by n matrix].
An m × n matrix is usually written as
a11 a12 … a1n
a21 a22 … a2n
a31 a32 … a3n
A = ……………………..
………………………
am1 am2 … amn
In a compact from the above matrix is represented by 𝐴 = 𝑎𝑖𝑗 , 𝑖 = 1,2, … , 𝑚, 𝑗 =
1,2, … , 𝑛 or simply by 𝑎𝑖𝑗 . We write the general element of the matrix and enclose
𝑚 ×𝑛

it in brackets of the type [ ] or of the type ( )

The numbers a11 , a12 etc. of this rectangular array are called the elements of the matrix.
The element aij belong to the ith row and the jth column and is sometimes called the
(i, j)th element of the matrix.
Thus in the elements aij the first suffix i will always denote the number of the row and
the second suffix 𝑗, the number of the column in which the element occurs.
In a matrix, the number of rows and columns need not be equal.
MATRIX MULTIPLICATION: The formal definition of matrix multiplication is as follows:

𝑎11 𝑎12 𝑎13 ⋯ 𝑎1𝑛

𝑎21 𝑎22 𝑎23 ⋯ 𝑎2𝑛

⋮
𝑎𝑚1 𝑎𝑚2 𝑎𝑚3 ⋯ 𝑎𝑚𝑛

𝑏11 𝑏12 𝑏13 ⋯ 𝑏1𝑝

𝑏21 𝑏22 𝑏23 ⋯ 𝑏2𝑝
×
⋮
𝑏𝑛1 𝑎𝑛2 𝑎𝑛3 ⋯ 𝑎𝑛𝑝

𝑛 𝑛 𝑛

𝑎1𝑖 𝑏𝑖𝑙 𝑎1𝑖 𝑏𝑖2 ⋯ 𝑎1𝑖 𝑏𝑖𝑝

𝑖=1 𝑖=1 𝑖=1
𝑛 𝑛 𝑛

𝑎2𝑖 𝑏𝑖𝑙 𝑎2𝑖 𝑏𝑖2 ⋯ 𝑎2𝑖 𝑏𝑖𝑝

𝑖=1 𝑖=1 𝑖=1
⋮
𝑛 𝑛 𝑛

𝑎𝑚𝑖 𝑏𝑖𝑙 𝑎𝑚𝑖 𝑏𝑖2 ⋯ 𝑎𝑚𝑖 𝑏𝑖𝑝

𝑖=1 𝑖=1 𝑖=1

The formula for matrix multiplication looks very complicated , but we can make more
sense of it by using the dot product of two vectors. The dot product of two vector is
formed by multiplying together each pair of corresponding components and then
adding the results of all these products.

A matrix can be thought of as either a vertical stack of row vectors:

𝑎11 ⋯ 𝑎1𝑛 𝑎1
⋮ = ⋮
𝑎𝑚1 ⋯ 𝑎𝑚𝑛 𝑎𝑚
 𝑎1 = 𝑎11 𝑎12, … , 𝑎𝑖𝑛

𝑎𝑚 = 𝑎𝑚1 𝑎𝑚2, … , 𝑎𝑚𝑛

Or as a horizontal stack of column vectors:

𝑏11 ⋯ 𝑏1𝑝
⋮ = 𝑏1 𝑏2, … , 𝑏𝑝
𝑏𝑛1 ⋯ 𝑏𝑛𝑝

𝑏11 𝑏1𝑝
𝑏1 = ⋮ ⋯ 𝑏𝑝 = ⋮
𝑏𝑛1 𝑏𝑛𝑝

For our purpose it is best to think of the left hand matrix 𝑨 as a collection of row
vectors, and the right hand matrix 𝑩 as a collection of column vectors. Then each
element in the matrix product AB can be found as a dot product of one row of 𝑨 with
one column of B:

𝑎1 ∙ 𝑏1 𝑎1 ∙ 𝑏2 𝑎1 ∙ 𝑏3 ⋯ 𝑎1 ∙ 𝑏𝑝
𝑎2 ∙ 𝑏1 𝑎2 ∙ 𝑏2 𝑎2 ∙ 𝑏3 ⋯ 𝑎2 ∙ 𝑏𝑝
𝐴𝐵 = ⋮
𝑎𝑚 ∙ 𝑏1 𝑎𝑚 ∙ 𝑏2 𝑎𝑚 ∙ 𝑏3 ⋯ 𝑎𝑚 ∙ 𝑏𝑝

The element in position 1,1 of the product matrix is the dot product of the first row of
𝑨 with the first column of B. In general, the element in position 𝑖, 𝑗 is formed by the
dot product of row 𝑖 in 𝐴 𝑎𝑛𝑑 𝑐𝑜𝑙𝑢𝑚𝑛 𝑗 𝑖𝑛 𝐵 . Examples of matrix multiplication are:

𝑎 𝑏 𝑒 𝑓 𝑎𝑒 + 𝑏𝑔 𝑎𝑓 + 𝑏𝑕
=
𝑐 𝑑 𝑔 𝑕 𝑐𝑒 + 𝑑𝑔 𝑐𝑓 + 𝑑𝑕

1 0 131211 38
11 12 13 100 1 232221 68
21 22 23 =
10000 2 333231 98
31 32 33

Matrix multiplication is a very valuable tool, making it much easier to write systems of
linear simultaneous equations. The three –equation system

𝑎1 𝑥 + 𝑏1 𝑦 + 𝑐1 𝑧 = 𝑑1
 𝑎2 𝑥 + 𝑏2 𝑦 + 𝑐2 𝑧 = 𝑑2

𝑎3 𝑥 + 𝑏3 𝑦 + 𝑐3 𝑧 = 𝑑3

Can be rewritten using matrix multiplication as

𝑎1 𝑏1 𝑐1 𝑥 𝑑1
𝑎2 𝑏2 𝑐2 𝑦 = 𝑑2
𝑎3 𝑏3 𝑐3 𝑧 𝑑3

n n
MATRIX OF A QUADRATIC FORM: If ∅ = i−1 j−1 a ij xi xj is a quadratic form in n
variabes x1 , x2 , … … , xn , then there exists a unique symmetrix matrix B of order n such
that ∅ = 𝑋 𝑇 𝐵𝑋 where 𝑋 = [x1 , x2 , … … , xn ]𝑇 . The symmetrix matrix B is called the
n n
matrix of the quadratic form i−1 j−1 a ij xi xj .

MATRIX OVER A FIELD 𝐅: A set of mn elements of a field F arranged in the form of a

rectangular array having m rows and n columns an m × n matrix over the field F.
If all the elements of a matrix belong to the field of real numbers, the matrix is said to be
real.
MATRIX POLYNOMIALS: An expression of the form

F λ = A° + A1 λ + A2 λ2 + ⋯ + Am−1 λm−1 + Am λm , where A0 , A1 , A2 , … , Am are all

square matrices of the same order, is called a Matric Polynomial of degree m provided.
Am is not a null matrix. The symbol λ is called indeterminate. According to this
definition of a matrix polynomial, each square matrix can be expressed as a matrix
polynomial with zero degree. For example, if A be any square matrix, we can write
A = λ°A.

MAXIMA: The maxima are the points where the value of a function is greater than it is at
the surrounding points.

MAXIMAL ASSIGNMENT PROBLEM: Sometimes, it may happen that instead of being

minimized, the problems utters to be maximized . In this case , we may follow the
slightly different but simple method for the solution, we can proceed as:

(a) We must first select the greatest element of our given assignment problem and
must subtract each element from this element to obtain a new matrix. The new matrix
obtained can now be solved by the usual method of assignment technique used for
minimal problems.
(b) One another method is also used for the conversion of maximization to
minimization. Place minus sign before each element of the given assignment matrix to
obtain a anew matrix. The new matrix thus can be treated as the usual assignment
minimal problem.

MAXIMAL ATLAS: The atlas containing all possible charts consistent with a given atlas is
called the maximal atlas (i.e. an equivalence class containing that given atlas (under the
already defined equivalence relation given in the previous paragraph)). Unlike an
ordinary atlas, the maximal atlas of a given manifold is unique. Though it is useful for
definitions, it is an abstract object and not used directly (e.g. in calculations).

MAXIMAL IDEAL: Let 𝑅 be a ring. A proper ideal 𝐼 of 𝑅 is said to be maximal if the only
ideals 𝐽 of 𝑅 satisfying 𝐼 ⊂ 𝐽 ⊂ 𝑅 are 𝐽 = 𝐼 and 𝐽 = 𝑅.
MAXIMUM DISTANCE SEPARABLE (MDS) CODE: A linear code with parameters [𝑛, 𝑘, 𝑑]
such that 𝑘 = 𝑛 − 𝑑 + 1 is called a maximum distance separable (MDS) code.
MAXIMUM FLOW PROBLEM:
 All flow through the directed, connected network originates from one node,
called the source; and, terminates at one other node, called the sink.
 All remaining nodes are called transshipment nodes.
 Directed arcs indicated direction of flow, and maximum amount of flow is given
by the arc capacity. At the source, all arcs point away; at the sink, all arcs point
into the node.
 Objective: maximize total amount of flow from source to sink. The amount
leaving the source is equal to the amount entering the sink.
MAXIMUM LIKELIHOOD DECODING: Let 𝐶 be a code of length 𝑛 over an alphabet 𝐴. The
maximum likelihood decoding rule states that every 𝑥 ∇ 𝐴𝑛 is decoded to 𝑐𝑥 ∇ 𝐶 when
𝑃[𝑥 𝑟𝑒𝑐𝑒𝑖𝑣𝑒𝑑 | 𝑐𝑥 𝑤𝑎𝑠 𝑠𝑒𝑛𝑡] = max 𝑃[𝑥 𝑟𝑒𝑐𝑒𝑖𝑣𝑒𝑑 | 𝑐 𝑤𝑎𝑠 𝑠𝑒𝑛𝑡]
𝑐 ∇𝐶

If there exist more than one 𝑐 with this maximum probability, then ⊥ is returned.
MAXIMUM LIKELIHOOD ESTIMATOR: A maximum likelihood estimator has this
property: if the true value of the unknown parameter is the same as the value of the
maximum likelihood estimator, then the probability of obtaining the sample that was
actually observed is maximizes.
MAXWELL’S EQUATIONS: Maxwell’s four equations govern electric and magnetic fields.
They were put together by James Clerk Maxwell in the 1870s on the basis of
experiments data. These equations can be used to establish the wave nature of light.

First, here are the equations for free space in integral form, assuming there is no change
in current over time.

Let E be an electric field (a three- dimensional vector field) Theses two equations apply:

1 𝐸 ∙ 𝑑𝐿 = 𝑜
𝑐𝑙𝑜𝑠𝑒𝑑 𝑝𝑎𝑡 𝑕

(That is, the line integral of the electric field over any closed path is zero)

𝑞𝑖𝑛𝑠𝑖𝑑𝑒
2 𝐸 ∙ 𝑑𝑆 =
𝑐𝑙𝑜𝑠𝑒𝑑 𝑝𝑎𝑡 𝑕 𝜀0

(That is, the surface integral of the electric field around any closed surface is equal to 𝑞,
the total charge inside the surface, divided by a constant known as 𝜀0

Let 𝑩 𝑏𝑒 a magnetic field (a three- dimensional vector field). Then the line integral
around a closed path depends on the current flowing through the interior of the path:

(3) ∫𝑝𝑎𝑡 𝑕 𝐿. 𝐵. 𝑑𝐿 = 𝜇𝑜 𝐼 𝑖𝑛𝑠𝑖𝑑𝑒

Where 𝐼 stands for the amount of electric current , and 𝜇𝑜 is a constant. The surface
integral of B over a closed surface is zero:

4 𝐵 ∙ 𝑑𝑆 = 0
𝑐𝑙𝑜𝑠𝑒𝑑 𝑝𝑎𝑡 𝑕

The four equations given above can also be written in alternate forms. Use Stroke’s
theorem to rewrite the two equations involving line integrals . Equation (1) becomes:

(5) ∆ ×𝐸 = 0

In words : the curl of the electric field is always zero.

Equation (3) becomes:

(6) ∆ × 𝐵 = 𝜇0 𝐽
Where 𝐽 , called the current density is defined by this integral:

𝐽 ∙ 𝑑𝑆 = 𝐼𝑖𝑛𝑠𝑖𝑑𝑒
𝑠𝑢𝑟𝑓𝑎𝑐𝑒 𝑆

By using the divergence theorem, the two equations involving surface integrals can be
rewritten. The left hand side of equation (2) is changed from

𝐸 ∙ 𝑑𝑆
𝑠𝑢𝑟𝑓𝑎𝑐𝑒 𝑆

Into

∆ ∙𝐸 𝑑𝑉
𝑖𝑛𝑡𝑒𝑟𝑖𝑜𝑟 𝑜𝑓 𝑆

The right-hand side of equation (2) is changed by defining 𝜌 over any volume is equal to
the total charge inside that volume:

𝑞𝑖𝑛𝑠𝑖𝑑𝑒 𝑆 = 𝜌𝑑𝑉
𝑖𝑛𝑡𝑒𝑟𝑖𝑜𝑟 𝑜𝑓 𝑆

We then have the equation

𝜌
∆ ∙𝐸 𝑑𝑉 = 𝑑𝑉
𝑖𝑛𝑡𝑒𝑟𝑖𝑜𝑟 𝑜𝑓 𝑆 𝑖𝑛𝑡𝑒𝑟𝑖𝑜𝑟 𝑜𝑓 𝑆 𝜀0

Since this equation must hold true for any arbitrary surface S, we can write the equation
in this form:

𝜌
(7) ∆∙𝐸 =
𝜀0

In words: the divergence of the electric field is proportional to the charge density.

Equation(4) becomes:

(8) ∆ ∙ 𝐵 = 0

MAXIMUM MODULUS PRINCIPLE: Suppose 𝑓 (𝑧) is analytic within and on a simple

closed contour 𝐶 and 𝑓(𝑧) is not constant. Then 𝑓(𝑧) reaches its maximum value on 𝐶
(and not inside 𝐶), that is to say, if 𝑀 is the maximum value of 𝑓(𝑧) on and within 𝐶,
then 𝑓(𝑧) < 𝑀 for every 𝑧 inside 𝐶.

MAX–MIN INEQUALITY: For any function

When equality holds one says that 𝐹, 𝑊, 𝑍 satisfies the strong max–min property

MAZUR–ULAM THEOREM: If 𝑉 and 𝑊 are normed spaces over R and the mapping

is a surjective isometry, then 𝑓 is affine.

m-DIMENSIONAL MANIFOLD: An m-dimensional manifold in 𝑅 𝑛 is a non-empty set

𝑆 ⊂ 𝑅 𝑛 satisfying the following property for each point 𝑝 ∇ 𝑆, there exists an open
neighborhood 𝑊 ⊂ 𝑅 𝑛 of 𝑝 and an m-dimensional embedded parametrized manifold
𝜍: 𝑈 → 𝑅 𝑛 with image 𝜍(𝑈) = 𝑆 ∩ 𝑊. The surrounding space 𝑅 𝑛 is said to be the
ambient space of the manifold.
m-DIMENSIONAL SMOOTH ATLAS: Let 𝑀 be a Hausdorff topological space, and let
𝑚 ≥ 0 be a fixed natural number. An m-dimensional smooth atlas of M is a collection
(𝑂𝑖 )𝑖∇𝐼 of open sets 𝑂𝑖 in 𝑀 such that 𝑀 = ⋃𝑖∇𝐼 𝑂𝑖 , together with a collection (𝑈𝑖 )𝑖∇𝐼 of
open sets in 𝑅 𝑚 and a collection of homeomorphisms, called charts, 𝜍𝑖 : 𝑈𝑖 → 𝑂𝑖 =
𝜍𝑖 (𝑈𝑖 ), with the following property of smooth transition on overlaps:
For each pair 𝑖, 𝑗 ∇ 𝐼 the map 𝜍𝑗 −1 ◦ 𝜍𝑖 is smooth from the open set 𝜍𝑖 −1 (𝑂𝑖 ∩ 𝑂𝑗 ) ⊂
𝑅 𝑚 to 𝑅 𝑚 .
MEAN PROPORTIONAL: The mean proportional is the geometric mean of two or more
numbers. If the 𝑛

numbers are 𝑥1 , 𝑥2 , 𝑥3 , … … . 𝑥𝑛 . the mean proportional is the nth root of their product.

𝑛
𝑥1 × 𝑥2 × 𝑥3 × … × 𝑥𝑛 .

MEASURABLE FUNCTION: A function 𝑓: 𝑋 → ℝ is measurable if 𝐸𝑐 𝑓 = {𝑥 ∇ 𝑋; 𝑓 𝑥 <

𝑐} is in 𝐿 for any 𝑐 ∇ ℝ. A complex-valued function is measurable if its real and
imaginary parts are measurable.
The following are equivalent:

1. A function f is measurable;

2. For any 𝑎 < 𝑏 the set 𝑓 −1 ((𝑎, 𝑏)) is measurable;

3. For any open set 𝑈 ∇ ℝ the set 𝑓 −1 (𝑈) is measurable.

Let 𝑓 be measurable and 𝑔 be continuous, then the composition 𝑔(𝑓(𝑥)) is measurable.

MEASURABLE SPACE: A measurable space is a set 𝐸 together with a collection 𝐵(𝐸) of

subsets of 𝐸 which is a sigma algebra. The elements of 𝐵(𝐸) are called measurable sets.
MEASURE OF OPEN INTERVALS: We define 𝜆(𝐼) = 𝑏 − 𝑎, where 𝐼 denotes the open
interval (𝑎, 𝑏). This is the beginning of a process that can, with some adjustments, be
applied to a variety of situations.
MEASURE OF OPEN SETS: Define 𝜆(𝐺) = 𝑋 𝜆(𝐼𝑘 ), where 𝐺 is an open set and {𝐼𝑘 } is the
sequence of component intervals of 𝐺. If one of the components is unbounded, we let
𝜆(𝐺) = ∞. [If 𝐺 ≠ ∅, then 𝐺 can be expressed as a finite or countably infinite disjoint
union of open intervals: 𝐺 = ⋃ 𝐼𝑘 . If 𝐺 = ∅, the empty set, define 𝜆(𝐺) = 0.] This
definition is a natural one; it conforms to our intuitive requirement that “the whole is
equal to the sum of the parts.”
MEASURE ON A RING: A measure is a map µ: 𝑅 → [0, ∞] defined on a ring (or
𝜍 −algebra) 𝑅, such that

1. µ(∅) = 0.

2. if (𝐴𝑛 ) is a finite subset of 𝑅, such that 𝐴𝑗 are pairwise disjoint (that is,
𝐴𝑛 ∩ 𝐴𝑚 = ∅ for 𝑛 ≠ 𝑚), then µ(⋃𝑛 𝐴𝑛 ) = 𝑛 µ(𝐴𝑛 ). This property is called
additivity of a measure. .

MEASURES OF CENTRAL TENDENCY: A measure of central tendency indicates a middle

or typical value of a group of numbers. Examples of measures of central tendency are
the mean, median or mode. Typically these three values are near each other, but not
always . For example, one very large value will significantly increase the value of the
mean. But it will not affect the median.

MEDIAN: (1) The median of a group of 𝑛 number is the number which divides the whole
data into two equal halves. For example the median of the set of numbers
1,2,3 𝑖𝑠 2; 𝑡𝑕𝑒 𝑚𝑒𝑑𝑖𝑎𝑛 𝑜𝑓 1,1,1,2,10,1516,20,100,110 is 15 . in order to determine
the median, the list should be placed in ascending or descending order. If there is an odd
number of items in the list , then the median is the element in the exact middle. If there
is an even number , then the medina is the average of the two numbers closest to the
middle.

(2) A median if a triangle is a line segment connecting one vertex to the midpoint of the
opposite side.

MERGELYAN'S THEOREM: Let 𝐾 be a compact subset of the complex plane 𝑪 such

that 𝑪\𝐾 is connected. Then, every continuous function 𝑓: 𝐾 𝑪, such that
the restriction 𝑓 to 𝑖𝑛𝑡(𝐾) is holomorphic, can be
approximated uniformly on 𝐾 with polynomials. Here, 𝑖𝑛𝑡(𝐾) denotes the interior of 𝐾.

MEROMORPHIC FUNCTION: A function which has poles as its only singularities in the
finite part of the plane is said to be meromorphic function.

MERSENNE PRIME: A Mersenne (also spelled Marsenne) prime is a specific type

of prime number. It must be reducible to the form 2𝑛 − 1, where 𝑛 is a prime number.
The term comes from the surname of a French monk who first defined it. The first few
known values of 𝑛 that produce Mersenne primes are where 𝑛 = 2, 𝑛 = 3, 𝑛 =
5, 𝑛 = 7, 𝑛 = 13, 𝑛 = 17, 𝑛 = 19, 𝑛 = 31, 𝑛 = 61, and 𝑛 = 89.

It takes the most powerful computer a long time to check a large number to determine if
it is prime, and an even longer time to determine if it is a Mersenne prime. For this
reason, Mersenne primes are of particular interest to developers of
strong encryption methods. In August 2008, Edson Smith, a system administrator at
UCLA, found the largest prime number known to that date. Smith had installed software
for the Great Internet Mersenne Prime Search (Gimps), a volunteer-based distributed
computing project. The number (which is a Mersenne prime) is 12,978,189 digits long.
It would take nearly two-and-a-half months to write out and, if printed, would stretch
out for 30 miles.

METABOLOMICS: Metabolomics is a term sometimes used to describe the emerging

science of measurement and analysis of metabolites, such as sugars and fats, in the cells
of organisms at specific times and under specific conditions. The field of metabolomics
overlaps with biology, chemistry, mathematics, and computer science. Metabolomics as
a discipline makes use of analytical processes such as spectroscopy, chromatography,
and multivariable analysis. Metabolomics allows scientists to measure physiological
effects and to monitor for adverse reactions to drugs. Metabolomics is of interest to
physicians because it may lead to improvements in the diagnosis and treatments of
human diseases.

𝜕𝑢 𝜕𝑢
METHOD OF CHARACTERISTICS: The first order wave equation 𝑐 𝜕𝑥 − 𝜕𝑡 = 0

describes the movement of a wave in one direction with no change of shape. The aim of
the method of characteristics is to solve the PDE by finding curves in the 𝑥 − 𝑡 plane
that reduce the equation to an ODE. In general, any curve in the 𝑥 − 𝑡 plane can be
expressed in parametric form by 𝑥 = 𝑥 𝑟 , 𝑡 = 𝑡 𝑟 , where the parameter, 𝑟, gives a
measure of the distance along the curve. The curve starts at the initial point, 𝑥 = 𝑥0 ,
𝑡 = 0, when 𝑟 = 0. Assuming that we can solve the resulting ODE means that 𝑢 is known
everywhere along this curve, i.e. along the curve picked out by the value of 𝑥0 . Another
choice for 𝑥0 gives another curve and the value of 𝑢 is determined along this curve. In
this manner, 𝑢 can be determined at any point in the 𝑥 − 𝑡 plane by choosing the curve,
defined by 𝑥0 , that passes through this point and taking the correct value of 𝑥0 , the
distance along the curve. Hence, we can evaluate 𝑢(𝑥, 𝑡).

Therefore, we have 𝑢 𝑥, 𝑡 = 𝑢(𝑥(𝑟), 𝑡(𝑟))and so 𝑢 is a function of . Hence, the

𝑑𝑢 𝑑𝑥 𝜕𝑢 𝑑𝑡 𝜕𝑢 𝑑𝑢
derivative of 𝑢 with respect to 𝑟 is = + 𝑑𝑟 . Consider = 0 provided the
𝑑𝑟 𝑑𝑟 𝜕𝑥 𝜕𝑡 𝑑𝑟
𝑑𝑥 𝑑𝑡
parametric representation of the curve satisfies = 𝑐1 and = 1 , which are the
𝑑𝑟 𝑑𝑟

characteristic curves. 𝑢 = 𝑐𝑜𝑛𝑠𝑡. along a characteristic curve but the constant may be
different on different characteristic curves. As 𝑥0 gives us a different characteristic
curve, this implies that 𝑢 = 𝐹(𝑥0 ).

Solving characteristic curves, we have 𝑡 = 𝑟, 𝑡 = 0 at 𝑟 = 0 and 𝑥 = 𝑐𝑟 + 𝑥0 = 𝑐𝑡 +

𝑥0 since 𝑥 = 𝑥0 𝑎𝑡 𝑟 = 0 ⇒ 𝑥0 = 𝑥 − 𝑐𝑡.

𝑥0 defines which characteristic curve we are working on, 𝑢 is constant on a

characteristic curve that depends on 𝑥0 . Hence, 𝑢 = 𝐹 𝑥0 = 𝐹(𝑥 − 𝑐𝑡), where the
arbitrary function, 𝐹, is determined by the initial condition.
METHOD OF EXHAUSTION: The method of exhaustion is calculating an area by
approximating it by the areas of a sequence of polygons.
METRIC SPACES: A metric space (𝑋, 𝑑) consists of a non-empty set 𝑋 with a distance
function 𝑑 defined on 𝑋 × 𝑋 such that

(i) 𝑑(𝑥, 𝑦) ≥ 0.
(ii) 𝑑(𝑥, 𝑦) = 0 𝑖𝑓𝑓 𝑥 = 𝑦.
(iii) 𝑑(𝑥, 𝑦) = 𝑑(𝑦, 𝑥).
(iv) 𝑑(𝑥, 𝑧) ≤ 𝑑(𝑥, 𝑦) + 𝑑(𝑦, 𝑧).

An open ball centred at 𝑥 ∇ 𝑋 is a set 𝐵(𝑥, 𝑟) = {𝑦 ∇ 𝑋 ∶ 𝑑(𝑥, 𝑦) < 𝑟}

METRIC TENSOR: A metric tensor defines how to measure distances along a path using
an integral based on a particular set of coordinated. For the simplest example, consider
this integral in two-dimensional Euclidian space using Cartesian coordinates. The
calculation is based on the differential distance 𝑑𝑆:

𝑑𝑆 2 = 𝑑𝑥12 + 𝑑𝑥22

Now write the expression like this:

𝑑𝑆 2 = 𝑔11 𝑑𝑥12 + 𝑔12 𝑑𝑥1 𝑑𝑥2 + 𝑔21 𝑑𝑥2 𝑑𝑥1 + 𝑔22 𝑑𝑥22

Where 𝑔 is matrix whose components are defied as:

𝑔11 = 1, 𝑔12 = 0,

𝑔21 = 0 𝑔22 = 1

The matrix 𝑔 is the metric (or the metric tensor). The components of 𝑔 determine how
to measure distances along curves in this particular space using these coordinates.

Written with summation notation:

2 2

𝑑𝑆 2 = 𝑔𝑖𝑗 𝑑𝑥𝑖 𝑑𝑥𝑗

𝑖=1 𝑗 =1
MEUSNIER’S THEOREM: Statement: If 𝜅 and 𝜅𝑛 are the curvatures of oblique and
normal sections through the same tangent line and 𝜃 be the angle between these
sections, then 𝜅𝑛 = 𝜅 𝑐𝑜𝑠𝜃.

z+z z−z
MILNE’S THOMSON’S METHOD: We have 𝑧 = 𝑥 + 𝑖𝑦 so that x = ,y= .
2 2i

𝑤 = 𝑓 𝑧 = 𝑢 + 𝑖𝑣 = 𝑢 𝑥, 𝑦 + 𝑖𝑣 (𝑥, 𝑦)

𝑧+𝑧 𝑧−𝑧 𝑧+𝑧 𝑧−𝑧

or 𝑓 𝑧 = 𝑢 , + 𝑖𝑣 ,
2 2𝑖 2 2𝑖

In fact, this relation is formal identity in two independent variables 𝑧 𝑎𝑛𝑑 𝑧

By setting 𝑥 = 𝑧, 𝑦 = 0, so that = 𝑧 , we obtain

𝑓 𝑧 = 𝑢 𝑧, 0 + 𝑖𝑣 (𝑧, 0)

We know that

𝑑𝑤 𝜕𝑤 𝜕𝑢 𝜕𝑣
𝑓′ 𝑧 = = = + 𝑖
𝑑𝑧 𝜕𝑥 𝜕𝑥 𝜕𝑥

𝜕𝑢 𝜕𝑢
= − 𝑖 𝜕𝑦 , by Cauchy-Riemann equations,
𝜕𝑥

𝜕𝑢
Taking 𝜕𝑥 = 𝜙1 𝑥, 𝑦 = 𝜙1 (𝑧, 0)

𝜕𝑢
= 𝜙2 𝑥, 𝑦 = 𝜙2 (𝑧, 0)
𝜕𝑦

We get 𝑓 ′ 𝑧 = 𝜙1 𝑧, 0 − 𝑖 𝜙2 (𝑧, 0)

Integeration yields the result

𝑓 𝑧 = 𝜙1 𝑧, 0 − 𝑖𝜙2 (𝑧, 0) 𝑑𝑧 + 𝑐,

Where 𝑐 is a constant. We can calculate 𝑓(𝑧) directly 𝑖𝑓 𝑢 is known

Similarly if 𝑣 (𝑥, 𝑦) is given, then it can be proved that

𝑓 𝑧 = Ψ1 𝑧, 0 + 𝑖Ψ2 (𝑧, 0) 𝑑𝑧 + 𝑐′
 𝜕𝑣 𝜕𝑣
Where Ψ1 = 𝜕𝑦 , Ψ2 = 𝜕𝑥

MINDING THEOREM (DIFFERENTIAL GEOMETRY): If two surface S and S’ have the same
constant curvatures, then they are locally isometric.

MINIMA: The minima are the points where the value of a function is less than it is at the
surrounding points.

MINIMAL POLYNOMIAL OF A MATRIX: The monic polynomial of lowest degree that

annihilates a matrix 𝐴 is called the minimal polynomial of 𝐴. Also if 𝑓 x is the minimal
polynomial matrix 𝐴. The minimal polynomial of a matrix is a divisor of every
polynomial that annihilates this matrix. The minimal polynomial of a matrix is a divisor
of the characteristic polynomial of that matrix.

MINIMAL SURFACE: If mean curvature of a surface is zero at all points, then the surface
is called a minimal surface.

MINIMAX AND MAXMIN PRINIPLES: If a player lists his worst possible outcomes of all
his potential strategy, then he will choose the best strategy among all these outcomes.

Such a principle is known as maxmin principle or optimal strategy

MINIMUM MODULUS PRINCIPLE: Suppose 𝑓(𝑧) is analytic within and on a closed

contour 𝐶 and let 𝑓 𝑧 ≠ 0 inside 𝐶.

Suppose further that 𝑓(𝑧) is not constant. Then 𝑓(𝑧) attains its minimum value at a
point on the boundary of 𝐶, that is to say, if 𝑚 is the minimum value of 𝑓(𝑧) inside and
on 𝐶, then

𝑓(𝑧) inside and on 𝐶, then

𝑓(𝑧) > 𝑚 ∀ 𝑧 𝑖𝑛𝑠𝑖𝑑𝑒 𝐶.

MINIMUM- PRINCIPLE NECESSARY OPTIMALITY THEOREM: Let θ be a numerical

function, let 𝑔 be an 𝑚- dimensional vector function, and let 𝑕 be an 𝑘- dimensional
vector function, all defined on some open set containing X 0 . Let x be a solution of

θ x = minx∇X θ x where 𝑥 ∇ 𝑋 = x: x ∇ X 0 , g x ≤ 0, h x = 0 .
Let θ and 𝑔 be differentiable at x, and let h have continuous first partial derivatives at x.
then there exists r ∇ R, r ∇ Rm , s ∇ Rk such that the following conditions are satisfied

[rθ ∆θ x + r ∆θ x + s ∆h x ](x − x) ≥ 0 for all x ∇ X ° where 𝑟g(x) = 0, (r0 , r) ≥ 0 and

(r0 , r, s) ≠ 0.

MINIMUM RATIO TEST (LINEAR PROGRAMMING): Choose strictly positive coefficients

in pivot column. Divide coefficients into right hand side of same row. Recognize row
with smallest ratio, this is the pivot row. Substitute variable of that row in entering basic
variable. Solve for new Basic feasible solution using elementary row operations.
MINIMUM SPANNING TREE: Given a graph G with weighted edges, a minimum spanning
tree is a spanning tree with minimum weight, where the weight of a spanning tree is the
sum of the weights of its edges. There may be more than one minimum spanning tree
for a graph, since it is the weight of the spanning tree that must be minimum.
MINIMUM WEIGHTED PATH LENGTH: Given a list of weights, 𝑊 = {𝑤1 , 𝑤2 , . . . , 𝑤𝑛 }, the
minimum weighted path length is the minimum of the weighted path length of all
extended binary trees that have n external nodes with weights taken from W. There
may be multiple possible trees that give this minimum path length, and quite often
finding this tree is more important than determining the path length.
MINKOWSKI: Herman Minkowski (1864 to 1909) developed the idea of four-
dimensional space –time , a concept used in relativity theory.

MINKOWSKI’S INEQUALITY: For 1 < 𝑝 < ∞, and 𝑛 ≥ 1, let u, 𝑣 ∇ 𝐾 𝑛 . Then

1 1
𝑛 𝑛 𝑝 𝑛 𝑞
𝑝 1 𝑝 𝑝 𝑞
( 𝑢𝑖 +𝑣𝑖 ) ≤ 𝑢𝑖 + 𝑣𝑖
𝑖=1 𝑖=1 𝑖=1

MINKOWSKI’S INEQUALITY FUNCTIONAL ANALYSIS): Let 𝑝 be a real number such that

𝑝 ≥ 1. Then for any

𝑓 = 𝑓 1 ,𝑓 2 ,… 𝑎𝑛𝑑 𝑔 = 𝑔 1 , 𝑔 2 , … .

belonging to 𝑙 𝑃 ,

∞ 1/𝑝 ∞ 1/𝑝 ∞ 1/𝑝

𝑝 𝑝 𝑝
𝑓 𝑘 +𝑔 𝑘 ≤ 𝑓 𝑘 + 𝑔 𝑘
𝑘=1 𝑘=1 𝑘=1
MIN-MAX THEOREM: Let 𝑨 be a 𝒏 × 𝒏 Hermitian matrix with eigenvalues 𝝀𝟏 ≥ . . . ≥
𝝀𝒌 ≥ . . . ≥ 𝝀𝒏 then

and

in particular,

and these bounds are attained when 𝑥 is an eigenvector of the appropriate eigenvalues.

Also note that the simpler formulation for the maximal eigenvalue 𝝀𝟏 is given by:

Similarly, the minimal eigenvalue 𝝀𝒏 is given by:

MINOR: Consider the determinant

a11 a12 a13

a
∆= 21 a22 a23 .
a31 a32 a33

If we leave the row and the column passing through the element aij , then the second
order determinant thus obtained is called the minor of the element aij and we denote it
by Mij . In this way we can get 9 minors corresponding to the 9 elements of ∆.

For example,

a12 a13
The minor of the element a21 = a a33 = M21,
32

a11 a13
The minor of the element a32 = a a23 = M32,
21

a22 a23
The minor of the element a11 = a a33 = M11, and so on.
32

In terms of the notation of minors, if we expand ∆ along the first row, then

∆= (−1)1+1 a11 M11 + −1 1+2

a12 M12 + −1 1+3
a13 M13
 = a11 M11 − a12 M12 + a13 M13

Similarly, if we expand ∆ along the second column, then

∆= −a12 M12 + a22 M22 − a32 M32 .

Thus we can express the determinant as a linear combination of the minors of the
elements of any row or any column.

MINOR ARC: A minor arc of a circle is an are with a measure less than 180.

MINOR AXIS: The minor axis of an ellipse is the line segment that passes through the
center of the ellipse that is perpendicular to the major axis.

MINUTE: A minute is a unit of measure for small angles equal to 1/60 of a degree,

MITTAG LEFFLER’S EXPANSION THEOREM: Suppose that the only singularities of 𝑓 𝑧 in

the finite part of the 𝑧 − 𝑝𝑙𝑎𝑛𝑒 are simple poles at 𝑎1 , 𝑎2 , … . 𝑎𝑛 arranged in the order of
increasing absolute values. Also suppose that

(i) Residues of 𝑓 𝑧 at 𝑎1 , 𝑎2, … , 𝑎𝑛 𝑏𝑒 𝑏1 , 𝑏2 … , 𝑏𝑛 ,

(ii) < 𝐶𝑛 > is a sequence of circles (or rectangles of square) of radii 𝑅𝑛 𝑜𝑟 𝑅𝑛 is
the minimum distance of 𝐶𝑛 , from the origin) 𝐶𝑛 encloses 𝑎1 , 𝑎2 ,…, 𝑎𝑛 and no
other poles. On the circle 𝐶𝑛 , 𝑓(𝑧) < 𝑀, where 𝑀 is independent of
𝑛 𝑎𝑛𝑑 𝑅𝑛 ⟶ ∞ as 𝑛 ⟶ ∞
Then for all values of 𝑧 except poles
∞
1 1
𝑓 𝑧 = 𝑓 0 + 𝑏𝑛 +
𝑧 − 𝑎𝑛 𝑎𝑛
𝑛=1

Number of poles and zeroes of a meromorphic function. Let 𝑓 z be analytic inside and
on a simple closed curve C except for a finite number of poles inside C, and let 𝑓 z ≠ 0
on C. Then

1 𝑓 z dz
2πi
∫c 𝑓 z
=N−P

Where N and 𝑃 are respectively the number the zeroes and the number of poles of 𝑓 z
inside 𝐶. A pole or zero of order 𝑛 is counted 𝑛 times.
MOBIUS TRANSFORMATION: A Mobius transformation is a bijection on the extended
complex plane 𝐶 ∪ {∞} given by
𝑎𝑧 + 𝑏 𝑑
;𝑧 ≠ − ,∞
𝑐𝑧 + 𝑑 𝑐
𝑎
𝑓 𝑧 = ;𝑧 = ∞
𝑐
𝑑
∞; 𝑧 = −
𝑐
where 𝑎, 𝑏, 𝑐, 𝑑 ∇ 𝐶 and 𝑎𝑑 − 𝑏𝑐 ≠ 0
It can be shown that the inverse, and composition of two mobius transformations are
similarly defined, and so the Mobius transformations form a group under composition.
MODE: The mode of a group of numbers is the number that occurs most frequently in
that group.

MODEL –I M/M/1 : ∞/𝑭𝑰𝑭𝑶 : This is a system with Poisson input, exponential waiting
time and Poisson output with single channel . Queue capacity of the system being
infinite with first in first out mode.

Notations: First M in the notation stands for Poisson input.

2nd M in the notation stands for Poisson output.

1 in the notation stands for number of channel.

∞ in the notation stands for infinite capacity of the system.

FIFO in the notation stands for first in first out.

MODIFIED MEAN VALUE THEOREM: Suppose 𝑓 ∶ 𝑅 𝑛 → 𝑅 𝑚 is differentiable with

||𝐷𝑓((1 − 𝑡)𝑎 + 𝑡𝑏)𝑕|| ≤ 𝑀||𝑕||

for all 𝑕 ∇ 𝑅 𝑛 and all 0 ≤ 𝑡 ≤ 1. Then

||𝑓(𝑏) − 𝑓(𝑎)|| ≤ 𝑀||𝑏 − 𝑎||.

MODIFIED PARSEVAL THEOREM: The map 𝑊: 𝐻 → 𝑙2 given by the formula

[𝑊𝑓](𝑛) = ⟨ 𝑓, 𝑒𝑛 ⟩ is an isometry for any orthonormal basis (𝑒𝑛 ).

MODULE: Let 𝑅 be a unital commutative ring. A set 𝑀 is said to be a module over 𝑅 (or
𝑅-module) if
 (i) given any 𝑥, 𝑦 ∇ 𝑀 and 𝑟 ∇ 𝑅, there are well-defined elements 𝑥 + 𝑦
and 𝑟𝑥 of 𝑀,
(ii) 𝑀 is an Abelian group with respect to the operation + of addition,
(iii) the identities 𝑟(𝑥 + 𝑦) = 𝑟𝑥 + 𝑟𝑦, (𝑟 + 𝑠)𝑥 = 𝑟𝑥 + 𝑠𝑥, (𝑟𝑠)𝑥 =
𝑟(𝑠𝑥), 1𝑥 = 𝑥 are satisfied for all 𝑥, 𝑦 ∇ 𝑀 and 𝑟, 𝑠 ∇ 𝑅.

If 𝐾 is a field, then a 𝐾-module is by definition a vector space over 𝐾. Let (𝑀, +) be an

Abelian group, and let 𝑥 ∇ 𝑀. If 𝑛 is a positive integer then we define 𝑛𝑥 to be the sum
𝑥 + 𝑥 + · · · + 𝑥 of 𝑛 copies of 𝑥. If 𝑛 is a negative integer then we define 𝑛𝑥 =
−(|𝑛|𝑥), and we define 0𝑥 = 0. This enables us to regard any Abelian group as a
module over the ring 𝑍 of integers. Conversely, any module over 𝑍 is also an Abelian
group. Any unital commutative ring can be regarded as a module over itself in the
obvious fashion. Let 𝑅 be a unital commutative ring, and let 𝑀 be an 𝑅-module. A subset
𝐿 of 𝑀 is said to be a submodule of 𝑀 if 𝑥 + 𝑦 ∇ 𝐿 and 𝑟𝑥 ∇ 𝐿 for all 𝑥, 𝑦 ∇ 𝐿 and
𝑟 ∇ 𝑅. If 𝑀 is an 𝑅-module and 𝐿 is a submodule of 𝑀 then the quotient group 𝑀/𝐿 can
itself be regarded as an 𝑅-module, where 𝑟(𝐿 + 𝑥) ≡ 𝐿 + 𝑟𝑥 for all 𝐿 + 𝑥 ∇ 𝑀/𝐿
and 𝑟 ∇ 𝑅. The 𝑅-module 𝑀/𝐿 is referred to as the quotient of the module 𝑀 by the
submodule 𝐿. Note that a subset I of a unital commutative ring 𝑅 is a submodule of 𝑅 if
and only if 𝐼 is an ideal of 𝑅. Let 𝑀 and 𝑁 be modules over some unital commutative ring
𝑅. A function 𝜙: 𝑀 → 𝑁 is said to be a homomorphism of 𝑅-modules if 𝜙(𝑥 + 𝑦) =
𝜙(𝑥) + 𝜙(𝑦) and 𝜙(𝑟𝑥) = 𝑟𝜙(𝑥) for all 𝑥, 𝑦 ∇ 𝑀 and 𝑟 ∇ 𝑅. A homomorphism of 𝑅-
modules is said to be an isomorphism if it is invertible. The kernel 𝑘𝑒𝑟 𝜙 and
𝑖𝑚𝑎𝑔𝑒 𝜙(𝑀) of any homomorphism 𝜙: 𝑀 → 𝑁 are themselves 𝑅-modules. Moreover if
𝜙: 𝑀 → 𝑁 is a homomorphism of 𝑅-modules, and if 𝐿 is a submodule of 𝑀 satisfying
𝐿 ⊂ 𝑘𝑒𝑟 𝜙, then 𝜙 induces a homomorphism 𝜙: 𝑀/𝐿 → 𝑁. This induced
homomorphism is an isomorphism if and only if 𝐿 = 𝑘𝑒𝑟 𝜙 and 𝑁 = 𝜙(𝑀).

MODULE HOMOMORPHISM: If 𝑀 and 𝑁 are modules over a ring 𝑅 we say that

𝜑 ∶ 𝑀 → 𝑁 is a module homomorphism if 𝜑(𝑟1 𝑚1 + 𝑟2 𝑚2 ) = 𝑟1 𝜑(𝑚1 ) + 𝑟2 𝜑(𝑚2 ) for all
𝑟1 , 𝑟2 ∇ 𝑅 and 𝑚1 , 𝑚2 ∇ 𝑀. If 𝜑 is a bijection we say that it is a module isomorphism
and that 𝑀 and 𝑁 are isomorphic.

MODULES OVER A UNITAL COMMUTATIVE RING: Let 𝑅 be a unital commutative ring. A

set 𝑀 is said to be a module over 𝑅 (or 𝑅-module) if
 (i) given any 𝑥, 𝑦 ∇ 𝑀 and 𝑟 ∇ 𝑅, there are well-defined elements 𝑥 + 𝑦
and 𝑟𝑥 of 𝑀,
(ii) 𝑀 is an Abelian group with respect to the operation + of addition,
(iii) the identities 𝑟(𝑥 + 𝑦) = 𝑟𝑥 + 𝑟𝑦, (𝑟 + 𝑠)𝑥 = 𝑟𝑥 + 𝑠𝑥, (𝑟𝑠)𝑥 =
𝑟(𝑠𝑥), 1𝑥 = 𝑥 are satisfied for all 𝑥, 𝑦 ∇ 𝑀 and 𝑟, 𝑠 ∇ 𝑅.

MODULI OF CONTINUITY: The modulus of continuity of a function 𝑓 ∇ 𝐶(𝑇) is defined

as

𝜔 𝑓, 𝑡 : = sup𝑕∇ 0,𝑡 sup𝑥∇𝑇 ∆1𝑕 𝑓, 𝑥 , ∆1𝑕 𝑓, 𝑥 = 𝑓(𝑥 + 𝑕) − 𝑓(𝑥).

Note that

(a) 𝜔(𝑓, 𝑡) is a continuous, non-negative and non-decreasing function of 𝑡;

(b) 𝜔(𝑓 + 𝑔, 𝑡) ≤ 𝜔(𝑓, 𝑡) + 𝜔(𝑔, 𝑡);

(c) 𝜔(𝑓, 𝑡) ≤ 2 𝑘𝑓𝑘;

(d) 𝜔(𝑓, 𝑛𝑡) ≤ 𝑛 𝜔(𝑓, 𝑡), 𝑛 ∇ 𝑁,

(e) 𝜔(𝑓, 𝑡) ≤ 𝑡 𝑓 ′ ;

(f) 𝜔(𝑓, 𝜆𝑡) ≤ (𝜆 + 1) 𝜔(𝑓, 𝑡), 𝜆 ∇ 𝑅.

MODULI OF SMOOTHNESS: The k-th modulus of smoothness of 𝑓 ∇ 𝐶(𝑇) is defined by

𝜔𝑘 𝑓, 𝑡 : = sup sup ∆𝑘𝑕 𝑓, 𝑥 .

𝑕∇ 0,𝑡 𝑥∇𝑇

Note that

(a) 𝜔𝑘 𝑓, 𝑡 is a continuous, non-negative and non-decreasing function of 𝑡;

(b) 𝜔𝑘 𝑓, 𝑛𝑡 ≤ 𝑛 𝑘 𝜔𝑘 𝑓, 𝑡 , 𝑛 ∇ 𝑁 ;

(c) 𝜔𝑘 𝑓, 𝑡 ≤ 𝑡 𝑘 𝑓 (𝑘) ;

(d) 𝜔𝑘 (𝑓, 𝜆𝑡) ≤ (𝜆 + 1)𝑘 𝜔𝑘 (𝑓, 𝑡), 𝜆 ∇ 𝑅 .

MODULUS: (1) In division, the modulus is the same as the remainder.

(2) The modulus of a complex number is the same as its absolute value.
MODULUS OF A COMPLEX NUMBER: The ordinary Euclidean distance of (𝑥, 𝑦) to (0, 0) is
𝑥 2 + 𝑦 2 . We also call this number the modulus of the complex number 𝑧 = 𝑥 + 𝑖𝑦

and we write |𝑧| = 𝑥 2 + 𝑦 2 . Note that 𝑧 · 𝑧 = 𝑥 2 + 𝑦 2 = |𝑧|2 . (1.3) The distance

from 𝑧 to 𝑤 is |𝑧 − 𝑤|. We also have the easily verified formulas |𝑧𝑤| = |𝑧||𝑤| and
|𝑅𝑒 𝑧| ≤ |𝑧| and |𝐼𝑚 𝑧| ≤ |𝑧|.

MODULUS OF A VECTOR:- The non-negative number which is the measure of the

magnitude of a vector is called its modulus or module. Thus the length of the line
segment OP is the modulus of 𝑂𝑃 . The modulus 𝛼 of a vector 𝑎 is sometimes written
as 𝛼 .

MODUS PONENES: Modus Ponens refers to an argument of the form:

Premise 1: If 𝐴, 𝑡𝑕𝑒𝑛 𝐵.

Premise 2: 𝐴 is true

Conclusion: 𝐵 𝑖s true.

MODUS TOLLENS: Modus tollens refers to an argument of the form:

Premise 1: If 𝐴, 𝑡𝑕𝑒𝑛 𝐵.

Premise 2: 𝐵 is true

Conclusion: 𝐴 is not true.

MOMENT OF A FORCE ABOUT A POINT (VECTOR MOMENT OR TORQUE): A force is

completely represented if we know its

(i) magnitude
(ii) Direction
(iii) Point of application.

It is known from the principle of transmissibility of force that the effect of a force
remains unaltered if the point of application is shifted to any point on the line of action
of the force. If we use a single free vector F to represent a force, it gives us only the
magnitude and direction of the force. To specify the line of action another vector is
necessary along with F. In the following section, we will define a vector called moment
vector, which specifies the line of action of the force. Hence a force is determined by

(i) the vector of the force and

(ii) the moment vector of the force about a specific point.

MOMENT OF A FORCE F ABOUT ANY STRAIGHT LINE THROUGH A GIVEN POINT O: The
moment q of a force F acting at a point P about a line L is a scalar given by

𝒒 = 𝒓 × 𝑭 . 𝒂.

Where 𝒂. is a unit vector in the direction of the line, and 𝑶𝑷 = 𝒓, 𝑶 being any point on
the line.

MONAD: A monad is:

1) A type of functor used in category theory in mathematics. Category theory describes

patterns in mathematical functions.

2) A kind of constructor used in input/output (I/O) operations without using language

features. In general, however, monads are useful whenever a programmer wants to
perform a purely functional computation separate from a related computation
performed apart from it.

3) A symbol used by ancient Greek philosophers, including Plato, Pythagoras and

Aristotle, to describe God or the totality of all beings. Metaphysical and theological
theory describes "monism" as the concept of "one essence."

MONIC POLYNOMIAL: A polynomial in 𝑥 in which the coefficient of the highest power of

𝑥 is unity is called a monic polynomial . The coefficient of the highest power of 𝑥 is also
called the leading coefficient of the polynomial. Thus x 3 − 2x 2 + 5x + 5 is a monic
polynomial of degree 3 over the field of real numbers. In this polynomial the leading
coefficient is 1.

MONGE’ THEOREM DIFFERENTIAL GEOMETRY): A necessary and sufficient condition

that a curve on a surface be a line of curvature is that the surface normals along the
curve form a developable.
MONOID: A monoid is a semigroup 𝐺 which contains an identity element; that is, there
exists an element 𝑒 ∇ 𝐺 such that 𝑒 · 𝑎 = 𝑎 · 𝑒 = 𝑎 for all 𝑎 ∇ 𝐺.
MONOMIAL: A monomial is an algebraic expression that does not involve any additions
4
or subtractions. For example, 4 × 3, 𝑎2 𝑏 3 and 3 𝜋𝑟 3 are all monomials.

MONOTONE CLASS THEOREM: The smallest monotone class containing an algebra of

sets 𝐺 is precisely the smallest 𝜍-algebra containing 𝐺.
MONOTONE CONVERGENCE THEOREM: Every bounded monotonic sequence converges.
MONOTONIC OPERATOR: For a poset 𝑋, an operator 𝑇 is a monotonic operator if for all
𝑥, 𝑦 ∇ 𝑋, 𝑥 ≤ 𝑦 implies 𝑇(𝑥) ≤ 𝑇(𝑦).
MONTE-CARLO METHOD: The basis of Monte Carlo technique is random sampling of a
variable’s possible values. For this technique, some random numbers are required
which may be converted into random verities whose behavior is known from last
experience. The main steps of Monte Carlo method are as follows:

Step (i) In order to have a general idea of the system, we first draw a flow diagram
of the system.

Step (ii) Then, we take correct sample observations to select some suitable model
for the system. In this step we compute the probability distributions for the variable of
our interest.

Step(iii) We then convert the probability distributions to a cumulative distribution

function.

Step (iv) A sequence of random number is now selected with the help of random
number tables.

Step (v) Next, we determine the sequence of values of variables of interest with
the sequence of random numbers obtained in (iv).

Step (vi) Finally, we construct some standard mathematical function to the values
obtained in (v).

Advantages of Monte-Carlo method

 (i) These are helpful in finding solutions of complicated mathematical
expressions which is not possible otherwise.
(ii) By these methods, difficulties of trail and error experimentation are
avoided.

Disadvantages of Monte-Carlo method

(i) These are costly way of getting a solution of any problem.

(ii) These methods do not provide optimal answers of the problems. The
answers are good only when the size of the sample is sufficiently large.

MONTE CARLO SIMULATION: A Monte Carlo simulation uses a random number

generator to model series of events. This method is used when it is uncertain whether
or not a particular event will occur, but the probability of occurrence can be estimated.

MONTEL’S THEOREM COMPLEX ANALYSIS): Let 𝑓(𝑧) be analytic function of 𝑧, regular

in the half strip 𝑆 defined by 𝑎 < 𝑥 < 𝑏, 𝑦 > 0. If 𝑓(𝑧) is bounded in 𝑆, and
lim𝑦 →∞ 𝑓 𝑧 = 𝑙 𝑓𝑜𝑟 𝑎 < 𝑥 < 𝑏 then 𝑓 𝑧 ⟶ 𝑙 on every line 𝑥 = 𝑥0 𝑖𝑛 𝑆 and indeed
𝑓 𝑧 → 𝑙 uniformly for 𝑎 + 𝛿 ≤ 𝑥0 ≤ 𝑏 − 𝛿.

MONTEL'S THEOREM: A uniformly bounded family of holomorphic functions defined on

an open subset of the complex numbers is normal.

MORDELL-WEIL THEOREM: If 𝐸 is an elliptic curve defined over a number field 𝐾, then

the group of points with coordinates in 𝐾 is a finitely generated abelian group.
MORERA’S THEOREM: Suppose 𝑓 is a continuous function in the open disc 𝐷 such that
for any triangle 𝑇 contained in 𝐷

𝑓(𝑧) 𝑑𝑧 = 0,
𝑇
then 𝑓 is holomorphic.
MULTICOLLINEARITY: The multicollinearity problem in multiple regression arises
when two or more independent variables are highly correlated . In that case, it is
difficult to determine the individual effects of the different variables. In the extreme
case where two independent variables are perfectly correlated, the multiple regression
calculation cannot be performed because it would involve dividing by zero.

MULTIINDEX: Let 𝑛 ∇ 𝑁. Then a element 𝛼 ∇ 𝑁 𝑛 is called a multiindex.

MULTILINEAR MAP: Let 𝑉 𝑘 = 𝑉 × · · · × 𝑉 be the Cartesian product of 𝑘 copies of 𝑉 . A
map 𝜙 from 𝑉 𝑘 to a vector space U is called multilinear if it is linear in each variable
separately (i.e. with the other variables held fixed).
MULTINOMIAL: A multinomial is the sum of two or more monomials. Each monomial is
called a term. For example , 𝑎2 𝑏 3 + 6 + 4𝑏 5 is a multinomial with three terms.

MULTIPLE REGRESSIONS: Suppose that a dependent variable 𝑌 depends on some

independent variables 𝑋1 , 𝑋2, 𝑎𝑛𝑑 𝑋3 according to the equitation:

𝑌 = 𝛽1 𝑋1 + 𝛽2 𝑋2 + 𝛽3 𝑋3 + 𝜀

Where 𝛽1 , 𝛽2 , 𝛽3 𝑎𝑛𝑑 𝛽4 are unknown coefficients and 𝜀 is an random variable called the
error term. The problem in multiple regression is to use observed values of the
𝑋 ′ 𝑠 𝑎𝑛𝑑 𝑦 to estimate the values of the 𝛽 ′ 𝑠

MULTIPLICAND: In the equation 𝑎𝑏 = 𝑐 , 𝑎 𝑎𝑛𝑑 𝑏 are the multiplicands.

MULTIPLICATION OF A MATRIX BY A SCALAR: Let A be any 𝑚 × 𝑛 matrix and 𝑘 any

complex number called scalar. The 𝑚 × 𝑛 matrix obtained by multiplying every
elements of the matrix. A by 𝑘 is called the scalar multiple of A by 𝑘 and is denoted by
𝑘A or Ak. Symbolically, if A = aij , then 𝑘A = Ak = kaij
m×n m×n

3 2 −1
For example if 𝑘 = 2 and 𝐴 =
4 −3 1 2×3,

2 × 3 2 × 2 2 × −1 6 4 −2
Then 2𝐴 = =
2 × 4 2 × −3 2 × 1 8 −6 2 2×3.

MULTIPLICATION OF TWO MATRICES: Let A = aij and B = bik m×n be two

m×n

matrices such that the number of columns in A is equal to the rows in 𝐵. Then the 𝑚 × 𝑝
matrix C = cik m×p such that

Cik = aij bik

i×1

[note that the summation is with respect to the repeated suffix]

is called the product of the matrices A and B in that order and we write 𝐶 = 𝐴𝐵.
In the product AB, the matrix A is called the pre-factor and the matrix 𝐵 is called the
post-factor. Also we say that the matrix A has been post- multiplied by the matrix 𝐵 and
the matrix 𝐵 has been pre-multiplied by the matrix A.

The product AB of two matrices A and B exists if and only if the number of columns in A
is equal to the number of rows in 𝐵. Two such matrices are said to be comformable for
multiplication. If A is an 𝑚 × 𝑛 matrix and 𝐵 is an 𝑛 × 𝑝 matrix, then AB is an 𝑚 × 𝑝
matrix. Further if A = aij and B = bik n×p , then AB = cik m×p where
m×n
n

Cik = aij bjk = ai1 b1k + ai2 b2k + ⋯ + ain bnk

j=1

i.e., the (i, k)th element cik of the matrix AB is obtained by multiplying the corresponding
elements of the ith row of A and the k th column of B and then adding the products.

 If the product AB exists, then it is not necessary that the product BA will also
exist. For example if A is a 4 × 5 matrix and B is a 5 × 3 matrix, then the product
AB exists while the product BA does not exist.
 Matrix multiplication is associative if conformability is assured; i.e., A BC =
AB C if A, B, C are m × n, 𝑛 × 𝑝, 𝑝 × 𝑞 matrices respectively.
 Multiplication of matrices is distributive with respect to addition of matrices i.e.,

A B + C = AB + AC

where A, B, C are any three, m × n, 𝑛 × 𝑝, 𝑛 × 𝑝 matrices respectively.

 The multiplicaltion of matrices is not always commutative.
 If A be any m × n matrix and On,p be an 𝑛 × 𝑝 null matrix. Then AOn,p =
Om,p where Om,p is an m × p null matrix.
 The equation 𝐴𝐵 = 0 does not necessarily imply that at least one of the matrices
𝐴 and 𝐵 must be a zero matrix.
 The product of two matrices can be a zero matrix while neither of them is a zero
matrix.
 In the case of matrix multiplication if 𝐴𝐵 = 0, then it does not necessarily imply
that 𝐵𝐴 = 0.
 If 𝐴 be an m × n matrix, 𝐼𝑛 denotes the n- rowed unit matrix, it can be easily seen
that AIn = A = In A.
MULTIPLICATION PRINCIPLE: If two choices are to be made, one from a list of m
possibilities and the second from a list of n possibilities , and any choice from the first
list can be combined with any choice form the second list, then the fundamental
principle of counting says that there mn total ways of making the choices.

MULTIPLICATIVE IDENTITY: The number 1 is the multiplicative identity , because

1 × 𝑎 = 𝑎 𝑓𝑜𝑟 𝑎𝑙𝑙 𝑎.

MULTIPLICATIVE INVERSE: The multiplicative inverse of a number 𝑎 (written as 1/𝑎

𝑜𝑟 𝑎−1 ) is the number that when multiplied by , gives a result of 1:

1
𝑎× =1
𝑎

1
The multiplicative inverse is also called the reciprocal. For example, 2 is the reciprocal of

2. There exists a multiplicative inverse for every real number except zero.

MULTIPLY-PERIODIC FUNCTION: A periodic function which has more than one

fundamental period is said to be Multiply-periodic.
N

NAGATA–SMIRNOV METRIZATION THEOREM: A topological space is metrizable if and

only if it is regular, Hausdorff and has a ς-locally finite base. A ς-locally finite base is a
base which is a union of countably many locally finite collections of open sets.

NAKAYAMA’S LEMMA: Let 𝑀 be a finitely-generated module over a unital commutative

ring 𝑅, and let 𝐽 be an ideal of 𝑅. Suppose that 𝐽𝑀 = 𝑀. Then there exists an element a
of 𝐽 with the property that 𝑎𝑚 = 𝑚 for all 𝑚 ∇ 𝑀.

NAPIER: John Napier (1550 to 1617) was a Scottish mathematician who developed the
concept of logarithms.

NATURAL LOGARITHM: The natural logarithm of a positive number 𝑥 𝑤𝑟𝑖𝑡𝑡𝑒𝑛 𝑎𝑠 𝐼𝑛 𝑥

is the logarithm of 𝑥 𝑡𝑜 𝑡𝑕𝑒 𝑏𝑎𝑠𝑒, 𝑤𝑕𝑒𝑟𝑒 𝑒 = 2.71828 …The natural logarithm function
can also be defined by the definite integral

𝑥
𝐼𝑛 𝑥 = 𝑡1 𝑑𝑡
𝑡

NATURAL NUMBERS: The natural numbers are the set of numbers 1,2,3,4,5,6,7,8 … . .
This set of numbers is also called the counting numbers, since they are the numbers
used to count something. They can also be called the positive integers.
NEAREST POINT THEOREM: Let 𝐾 be a non-empty convex closed subset of a Hilbert
space 𝐻. For any point 𝑥 ∇ 𝐻 there is the unique point 𝑦 ∇ 𝐾 nearest to 𝑥.

 In ℝ2 with either norm ||·||1 or ||·||∞ , the nearest point could be non-unique.
 Let 𝑀 be a subspace of a Hilbert space 𝐻 and a point 𝑥 ∇ 𝐻 be fixed.
Then 𝑧 ∇ 𝑀 is the nearest point to 𝑥 if and only if 𝑥 − 𝑧 is orthogonal to any
vector in 𝑀.
 The minimal value of ||𝑥 − 𝑦|| for 𝑦 ∇ 𝐻1 is achieved when 𝑦 =
𝑛
𝑖=1⟨ 𝑥, 𝑒𝑖 ⟩ 𝑒𝑖 .

NEAR OPERATOR: Let 𝑋 be a set and 𝑌 a Banach space. Let 𝐴, 𝐵 be two operators from 𝑋
to 𝑌 . We say that 𝐴 is near 𝐵 if and only if there exist two constants 𝛼 > 0 and
𝑘 ∇ (0, 1) such that, for each 𝑥0 , 𝑦0 ∇ 𝑋, one has
𝐵(𝑥0 ) − 𝐵(𝑦0 ) − 𝛼(𝐴(𝑥0 ) − 𝐴(𝑦0 )) ≤ 𝑘 𝐵(𝑥0 ) − 𝐵(𝑦0 )

NECESSARY: In the statement 𝑝 → 𝑞. 𝑞 is a necessary condition for 𝑝 to be true. For

example, having four 90 angles is a necessary condition for a quadrilateral to be a
square (but it is not a sufficient condition).

NECESSARY AND SUFFICIENT CONDITION FOR R-INTEGRABILITY: A necessary and

sufficient condition for R-integrability of a bounded function 𝑓: [𝑎, 𝑏] → 𝑅 over [𝑎, 𝑏] is
that for every ∇> 0, there exists a partition 𝑃 of [𝑎, 𝑏] such that for 𝑃 and for every
refinement of 𝑃
0 ≤ 𝑈 𝑃, 𝑓 − 𝐿 𝑃, 𝑓 <∇.
NEGATION: The negation of a statement 𝑝 is the statement NOT 𝑝.

NEGATIVE: A negative number is any real number less than zero. The negative of any
number 𝑎 (written as – 𝛼) is defined by this equation : 𝑎 + −𝑎 = 0

NEGATIVE CORRELATION: A negative correlation is a relationship between

two variables such that as the value of one variable increases, the other decreases.
Correlation is expressed on a range from +1 to -1, known as the correlation coefficent.
Values below zero express negative correlation. A perfect negative correlation has a
coefficient of -1, indicating that an increase in one variable reliably predicts a decrease
in the other one. A perfect positive correlation, which has a coefficient of +1, indicates
that an increase or decrease in one variable always predicts the same directional change
for the second variable. Lower degrees of correlation are expressed by non-zero
coefficents between +1 and -1. Zero indicates a lack of correlation: There is no tendency
for the variables to fluctuate in tandem either positively or negatively.

NEGATIVE ORIENTATION: The negative orientation (counterclockwise) is the one that

is given by the standard parametrization 𝑧(𝑡) = 𝑧0 + 𝑟𝑒 −𝑖𝑡 , where 𝑡 ∇ 0, 2𝜋 .
NEGATIVE VECTOR: The vector which has the same modulus as the vector 𝜶 but
opposite direction, is called the negative of 𝜶.

The negative of a is represented by – . Thus if 𝐴𝐵 = 𝜶, then 𝐵𝐴= − 𝛼.

NEIGHBOURHOOD: Topology starts with a definition of open sets which are often and
interchangeably called neighborhoods. A neighbourhood of 𝑥 ∇ 𝑋 is a subset of the
topological space which contains an open set 𝑈 with 𝑥 ∇ 𝑈.
By definition, the empty set is open and so are finite intersections and arbitrary unions
of open sets. A set is a topological space if some of its subsets are declared to be open
subject to these conditions. An open set is a neighborhood of all its points. A
function 𝑓: 𝐴 → 𝐵 from one topological space into another is continuous at a point
𝑎 ∇ 𝐴 if for every neighborhood V of 𝑓(𝑎) there exists a neighborhood U of a such that
𝑓(𝑈) ⊂ 𝑉. This is equivalent to saying that inverse images of sets open in 𝐵 are open in
𝐴. Complements of open sets are closed by definition.
NESBITT’S INEQUALITY: Nesbitt’s inequality says, that for positive real 𝑎, 𝑏 and 𝑐 we
have:
𝑎 𝑏 𝑐 3
+ + ≥
𝑏+𝑐 𝑎+𝑐 𝑎+𝑏 2
NEUMANN’S INTEGRAL: If n is zero or a positive integer, then

+1
1 Pn (x)
Q1 y = 2 dx
−1 y−x

where y > 1 𝑎𝑛𝑑 − 1 ≤ x ≤ 1.

NEWTON: Sir Isaac Newton (1643 to 1727) was an English mathematician and scientist,
who developed the theory of gravitation and the laws of motion, designed a reflecting
telescope using a paraboloid mirror, used a prism to split white light into component
colors, and was one of the inventors of calculus.
NEWTON METHOD: Newton method provides a way to estimate the places where
complicated functions cross the x- axis. First, make a guess 𝑥1 that seems reasonably
close to the true value. Then approximate the curve by its tangent line to estimate a new
value 𝑥2 from the equation

𝑓 𝑥1
𝑥2 = 𝑥1 −
𝑓′ 𝑥1

where 𝑓′ 𝑥1 is the derivative of the function 𝑓 at the point 𝑥1 . The process is iterative:
that is it can be repeated as often as we like. This mean that we can get as close to the
true value as we wish.

NILPOTENT MATRIX: A non-zero matrix 𝐴 is said to be nilpotent, if the some positive

integer 𝑟, 𝐴′ = 0. A non-zero nilpotent matrix cannot be similar to a diagonal matrix.

NOETHER: Emmy Noether (1182 to 1935) was a German mathematician who

contributed to abstract algebra.

NOETHERIAN MODULES: Let 𝑅 be a unital commutative ring. An 𝑅-module 𝑀 is said to

be Noetherian if every submodule of 𝑀 is finitely-generated.
NOETHERIAN RING: A unital commutative ring is said to be a Noetherian ring if every
ideal of the ring is finitely-generated. A Noetherian domain is a Noetherian ring that is
also an integral domain.
NOETHER’S THEOREM: Consider a one-parameter family of maps 𝑞𝑖 𝑡 → 𝑄𝑖 𝑠, 𝑡 , 𝑠 ∇
𝑅 such that 𝑄𝑖 (0, 𝑡) = 𝑞𝑖 (𝑡). Then this transformation is said to be a continuous
symmetry of the Lagrangian 𝐿 if 𝜕 /𝜕𝑠 𝐿(𝑄𝑖 (𝑠, 𝑡), 𝑄𝑖 (𝑠, 𝑡), 𝑡) = 0. Noether’s theorem
states that for each such symmetry there exists a conserved quantity.
NOMOGRAMS: Nomograms are charts in which we can easily read off the corresponding
value 𝑢𝑛 , from given values 𝑢1 , 𝑢2 , − − −𝑢𝑛−1 when there is a relation 𝐹 𝑢1 , 𝑢2 , − −
−𝑢𝑛 = 0 among 𝑛 real variables 𝑢1 , 𝑢2 , − − −𝑢𝑛−1 . The construction of nomograms
has been thoroughly investigated by M. d’ocagne
NON-SINGULAR AND SINGULAR MATRICES: A square matrix A is said to be non-singular
ror singular according as A ≠ 0 or A = 0. Thus the necessary and sufficient condition
for a matraix to be invertible is that it is non-singular. If A, B be two 𝑛 -rowed non-
singular matrices, then AB is also non-singular and (AB)−1 = B−1 A−1 ,
If A be an 𝑛 × 𝑛 non-singular matrix, then A′ −1
= A−1 ′

−1
If A be an 𝑛 × 𝑛 non-singular matrix, then A′ θ
= Aθ .

NON-TRIVIAL CLOSED TRAIL: Let (𝑉, 𝐸) be a graph. A circuit in the graph is a non-trivial
closed trail in the graph.

NORM: Let 𝑆 be a linear space. The norm is an assignment of a non-negative real

number to every vector 𝑓 of 𝑆, written as 𝑓 , in such a way that

(i) 𝑓 =0 if and only if 𝑓 = 0

(so that 𝑓 >0, if and only if 𝑓 ≠ 0),
(ii) 𝑎𝑓 = 𝑎 𝑓 for any scalar a,
(iii) 𝑓+𝑔 ≤ 𝑓 + 𝑔 (the triangle inequality)

If a norm is defined on a linear space 𝑆, we say that it is normed linear space.

NORMAL: In mathematics the word “normal” means “perpendicular “. A line is normal to

a curve if it is perpendicular to a tangent line o that curve at the point where it
intersects the curve. Two vectors are normal to each other if their dot product is zero.

NORMAL ANGLE (DIFFERENTIAL GEOMETRY): The angle between the principal

normal and the surface normal is known as normal angle and is denoted by 𝜛. Hence
𝑁 . 𝑏 = sin 𝜛, 𝑁. 𝑛 = 𝑐𝑜𝑠𝜛.

NORMAL DISTRIBUTION: A continuous random variable X has a normal distribution if

its density function is

1 2/2𝑎𝑟 2
𝑓 𝑥 = 𝑒 − 𝑥−𝑝
𝜍 2𝜋

The mean (or expectation) of 𝑋 is 𝜇 and its variance is 𝜍 2 . If 𝜇 = 0 𝑎𝑛𝑑 𝜍 = 1 , 𝑡𝑕𝑒𝑛 𝑋 is

said to have the standard normal distribution , which has the density function

1 2/2
𝑓 𝑥 = 𝑒 −𝑥
2𝜋

The central limit theorem is one important application of the normal distribution . The
central limit theorem states that if 𝑋1 , 𝑋2 , … 𝑋𝑛 are independent, identically distrusted
random variables each with mean 𝜇 and variance 𝜍 2 , then in the limit that n goes to
infinity,

𝑆𝑛 = 𝑋1 + 𝑋1 + 𝑋3 + ⋯ + 𝑛

Will have a normal distribution with mean 𝑛𝜇 and variance 𝑛𝜍 2 . The reason that this
theorem is so remarkable is that it is completely general . It says that, no matter how 𝑋
is distributed , if you add up enough measurements, the sum of the 𝑋 ′ 𝑠 will have a
normal distribution. A normal distribution is an arrangement of a data set in which
most values cluster in the middle of the range and the rest taper off symmetrically
toward either extreme.

Height is one simple example of something that follows a normal distribution pattern:
Most people are of average height, the numbers of people that are taller and shorter
than average are fairly equal and a very small (and still roughly equivalent) number of
people are either extremely tall or extremely short.

Here’s an example of a normal distribution curve:

A graphical representation of a normal distribution is sometimes called a bell curve

because of its flared shape. The precise shape can vary according to the distribution of
the population but the peak is always in the middle and the curve is always
symmetrical. In a normal distribution, the mean, mode and median are all the same.

Normal distribution curves are sometimes designed with a histogram inside the curve.
The graphs are commonly used in mathematics, statistics and corporate data analytics.

NORMAL FUZZY SET: A fuzzy subset 𝐴 of a classical set 𝑋 is called normal if there exists
an 𝑥 ∇ 𝑋 such that 𝐴(𝑥) = 1. Otherwise A is subnormal.
NORMALISED FUNCTIONS OF BOUNDED VARIATION: The function 𝑔 𝑡 ∇ 𝐵𝑉 𝑎, 𝑏 is
said to be normalized function of bounded variation if 𝑔 𝑧 = 0 and 𝑔 is continuous for
all 𝑡 from the right, i.e. for all 𝑡 ∇ 𝑎, 𝑏 , 𝑔 𝑡 + 0 = 𝑔 𝑡 .

NORMALISED REPRESENTATION OF THE EQUIVALENCE CLASS OF 𝒇: Let 𝑓 𝑡 ∇

𝐵𝑉[𝑎, 𝑏] and let the function 𝑓(𝑡) be defined as follows:

𝑓 𝑡 = 0, at 𝑡 = 𝑎,

𝑓 𝑡 = 𝑓 𝑡 + 0 − 𝑓 𝑎 , at any interior point 𝑡 ∇ 𝑎, 𝑏 ,

and 𝑓 𝑡 = 𝑓 𝑏 − 𝑓, 𝑎𝑡 𝑡 = 𝑏

Then 𝑓(𝑡) is called a normalized representative of the equivalence class of 𝑓.

NORMALISED VECTORS: A vector is said to be normalized it its norm is 1. Any non-null

vector 𝑓 can be normalized by taking 𝑓 / 𝑓 . A normalized vector is also called a unit
vector.

NORMAL MATRIX: A matrix 𝐴 is said to be normal if 𝐴𝐴° = 𝐴°𝐴.

 The Hermitian, real symmetric, unitary, real orthogonal, skew- Hermitian, and
real skew- symmetric matrices are normal.
 Any diagonal matrix over the complex field is normal.
 Every matrix unitarily similar to a normal matrix is normal.
 If 𝑋 is an eigenvector of a normal matrix 𝐴 corresponding to an eigenvalue 𝜆,
then 𝑋 is also an eigenvector of 𝐴°; the corresponding eigenvalue being 𝜆.
 Characteristic vectors corresponding to distinct characteristic values of a normal
matrix are orthogonal.
 A triangular matrix is normal if and only if it is diagonal.

NORMAL OPERATOR: A normal operator 𝑇 is one for which 𝑇 ∗ 𝑇 = 𝑇𝑇 ∗ .

Note that

1. Any self-adjoint operator T is normal, since 𝑇 ∗ =𝑇.

2. Any unitary operator 𝑈 is normal, since 𝑈 ∗ 𝑈 = 𝑈𝑈 ∗ = 𝐼..

 3. Any diagonal operator 𝐷 is normal , since 𝐷 𝑒𝑘 = 𝜆𝑘 𝑒𝑘 , 𝐷 ∗ 𝑒𝑘 = 𝜆𝑘 𝑒𝑘 , and
𝐷𝐷∗ 𝑒𝑘 = 𝐷∗ 𝐷 𝑒𝑘 = | 𝜆𝑘 |2 𝑒𝑘 .

4. The shift operator S is not normal.

NORMAL STRAIN: It is defined as the ratio of the change is length to the original length
of a straight line element. It is also known as extension or compression.

NORMED SPACE (𝑳𝒑 (µ) / ∼, || · ||𝒑): We write 𝐿𝑝 (µ) for the normed space (𝐿𝑝 (µ) / ∼,
|| · ||𝑝 ). We will use notation and continue to write members of 𝐿𝑝 (µ) as functions.
Really they are equivalence classes, and so care must be taken when dealing with 𝐿𝑝 (µ).
For example, if 𝑓 ∇ 𝐿𝑝 (µ), it does not make sense to talk about the value of 𝑓 at a point.

Let (𝑓𝑛 ) be a Cauchy sequence in 𝐿𝑝 (µ). There exists 𝑓 ∇ 𝐿𝑝 µ with ||𝑓𝑛 − 𝑓||𝑝 → 0. In
fact, we can find a subsequence (𝑛𝑘 ) such that 𝑓𝑛 𝑘 → 𝑓 pointwise, almost everywhere.

𝐿𝑝 µ is a Banach space.

Let (𝑋, 𝐿, µ) be a measure space, and let 1 ≤ 𝑝 < ∞. We can define a map 𝛷: 𝐿𝑞 (µ)→
𝐿𝑝 (µ)* by setting 𝛷(𝑓) = 𝐹, for 𝑓 ∇ 𝐿𝑞 (µ), 1/𝑝 + 1/𝑞 = 1, where

𝐹: 𝐿𝑝 µ → 𝐾, 𝐹 𝑔 = ∫𝑋 𝑓𝑔𝑑𝜇 𝑔 ∇ 𝐿𝑝 µ

Let (𝑋, 𝐿, µ) be a finite measure space, let 1 ≤ 𝑝 < ∞, and let 𝐹 ∇ 𝐿𝑝 (µ)∗ . Then there
exists 𝑓 ∇ 𝐿𝑞 (µ), 1/𝑝 + 1/𝑞 = 1 such that

𝐹 𝑔 = ∫𝑋 𝑓𝑔𝑑𝜇 𝑔 ∇ 𝐿𝑝 µ

For 1 < 𝑝 < ∞, we have that 𝐿𝑝 (µ)∗ = 𝐿𝑞 (µ) isometrically, under the identification of
the above results. Note that 𝐿∞ ∗ is not isomorphic to 𝐿1 , except finite-dimensional
situation. Moreover if µ is not a point measure 𝐿1 is not a dual to any Banach space.

Let (𝑋, 𝐿, µ) be a finite measure space, and let 1 ≤ 𝑝 < ∞. Then the collection of simple
functions is dense in 𝐿𝑝 (µ).

For 1 < 𝑝 < ∞, we have that 𝐶𝐾 ([0,1]) is dense in 𝐿𝑝 ([0,1]).

NORM OF A PARTITION: The norm of a partition 𝑃 is the greatest of the lengths of the
segments of the partition 𝑃 and it is denoted by 𝑃 .
NORM OF LINEAR OPERATOR: A norm of linear operator is defined:
 𝑇 = sup⁡
{ 𝑇 𝑥 : 𝑥 ≤ 1}. 𝑇 is a bounded linear operator if
||𝑇|| = 𝑠𝑢𝑝{||𝑇𝑥||: ||𝑥||} < ∞.

Note that

1. Kernel of T is a linear subspace of X and image of T is a linear subspace of Y.

2. ||𝑇𝑥|| ≤ ||𝑇|| · ||𝑥|| 𝑓𝑜𝑟 𝑎𝑙𝑙 𝑥 ∇ 𝑋.

Consider the following examples and determine kernel and images of the mentioned
operators.

I. On a normed space X define the zero operator to a space 𝑌 by 𝑍: 𝑥 → 0 for all

𝑥 ∇ 𝑋. Its norm is 0.

II. On a normed space X define the identity operator by 𝐼𝑋 : 𝑥 → 𝑥 for all 𝑥 ∇ 𝑋. Its
norm is 1.

III. On a normed space 𝑋 any linear functional define a linear operator from 𝑋 to ℂ,
its norm as operator is the same as functional.

IV. The set of operators from ℂn to ℂm is given by n× m matrices which acts on

vector by the matrix multiplication. All linear operators on finite-dimensional
spaces are bounded.

V. Let 𝑆(𝑥1 , 𝑥2 , … ) = (0, 𝑥1 , 𝑥2 , … ) be the right shift operator. Clearly ||𝑆𝑥|| =

||𝑥|| for all 𝑥, 𝑠𝑜 ||𝑆|| = 1.

Let 𝑇: 𝑋 → 𝑌 be a linear operator. The following conditions are equivalent:

1. 𝑇 is continuous on 𝑋;

2. 𝑇 is continuous at the point 0.

3. 𝑇 is a bounded linear operator.

 Let 𝐵(𝑋, 𝑌) be the space of bounded linear operators from 𝑋 and 𝑌 with the norm
defined above. If 𝑌 is complete, then 𝐵(𝑋, 𝑌) is a Banach space.
 Let 𝑇 ∇ 𝐵(𝑋, 𝑌) and 𝑆 ∇ 𝐵(𝑌, 𝑍), where 𝑋, 𝑌, 𝑎𝑛𝑑 𝑍 are normed spaces. Then
𝑆𝑇 ∇ 𝐵(𝑋, 𝑍) and ||𝑆𝑇|| ≤ ||𝑆||||𝑇||.
 Let 𝑇 ∇ 𝐵(𝑋, 𝑋) = 𝐵(𝑋), where 𝑋 is a normed space. Then for any 𝑛 ≥ 1, 𝑇𝑛 ∇
𝐵(𝑋) 𝑎𝑛𝑑 ||𝑇𝑛|| ≤ ||𝑇||𝑛.
Let 𝑇 ∇ 𝐵(𝑋, 𝑌). We say 𝑇 is an invertible operator if there exists 𝑆 ∇ 𝐵(𝑌, 𝑋) such that
𝑆𝑇 = 𝐼𝑋 𝑎𝑛𝑑 𝑇𝑆 = 𝐼𝑌 . Such an 𝑆 is called the inverse operator of 𝑇.

1. The zero operator is never invertible unless the pathological spaces 𝑋 = 𝑌 = {0}.

2. The identity operator 𝐼𝑋 is the inverse of itself.

3. A linear functional is not invertible unless it is non-zero and X is one

dimensional.

4. An operator ℂn→ ℂm is invertible if and only if m=n and corresponding square

matrix is non-singular, i.e. has non-zero determinant.

NORMS OF A BOUNDED LINEAR FUNCTIONAL: Let 𝑁 be a normed linear space and Φ

abounded linear functional defined on 𝑁. Since Φ is bounded, there exists a non-
negative real number 𝑘 such that

Φ𝑓 ≤ 𝑘 𝑓 , for all 𝑓 ∇ 𝑁.

The smallest number 𝑘, such that Φ(𝑓) ≤ 𝑘 𝑓 holds for all 𝑓, is called the norm of Φ,
and is denoted by Φ . Thus

Φ = inf. k: k ≤ 0, Φ (f) ≤ k f ,

For all 𝑓 𝜖 𝑁.

NORMS OF MAPPINGS: Let 𝑋 be a normed linear space (where the norm provides the
distance function). Then, if 𝑀 is a linear mapping from 𝑋 to 𝑈, its norm is defined to be
the number 𝑀 : = 𝑠𝑢𝑝 { 𝑀(𝑓) ∶ 𝑓 ∇ 𝑋, 𝑓 = 1}. It follows from linearity of the
norm that the inequality 𝑀(𝑓) ≤ 𝑀 𝑓 holds for every 𝑓 ∇ 𝑋.

NORTH –WEST CORNER RULE: This is the most widely used method for the solution of
transportation problem. The following steps must be followed to solve the problem:

(i) The first step is to check the total of both demand and supply
𝑚 𝑛
i.e. , 𝑖=1 𝑎𝑖 = 𝑗 =1 𝑏𝑗 … 1

if the above equation satisfies the problem is said to be a balanced problem, otherwise
it is called an unbalanced problem.
(ii) The next step is to allot the corner assignments. In this step we start from the top
most cell at the North West corner and allocate in that cell the maximum amount
possible.
If 𝑥11 be the allocation in that very cell, then:
𝑥11 = min 𝑎1 , 𝑏1
(iii) Now, we have certain conditions to proceed further:
(a) If 𝑏1 < 𝑎1 , 𝑡𝑕𝑒𝑛 𝑥11 = 𝑏1 .

It means a little quantity is still left in the first row. So, we have to move to the right side
for the 2nd cell in the same row, i.e. , for cell (1,2) and make the allocation there which
must be of the order:

𝑥12 = min. 𝑎1 − 𝑥11 , 𝑏2

(b) If 𝑏1 > 𝑎1 , 𝑡𝑕𝑒𝑛 𝑥11 = 𝑎1. it means a little quantity has been left in column 1.
Then, we move downwards to the cell (2, 1) and make next allocation of amount
𝑥21 = min. 𝑎1 , 𝑏1 − 𝑥11 in the particular cell.
(c) If however, 𝑏1 = 𝑎1 , then no quantity is left to be allocated in the same row or
column, i.e., 𝑥12 = 0 𝑜𝑟 𝑥21 = 0.
(iv) Repeat the steps from (i) to (iii) until the entire requirement is satisfied.

NOT: The word “NOT” is used in logic to indicate the negation of a statement . The
statement “NOT p” IS FALSE IF P IS TRUE, AND IT IS TRUE IF P IS FALSE. The operation
of NOT can be described by this truth table.

𝑝 𝑁𝑂𝑇 𝑃

𝑇 𝐹

𝐹 𝑇

The symbols – 𝑝 𝑜𝑟 ~ 𝑝 𝑜𝑟 𝑝 are used to represent NOT.

𝒏 −PERSON GAME: When number of players is 𝑛 where 𝑛 ≥ 2 ,then for 𝑛 = 2 , the game
is known as 2 person game and for 𝑛 > 2 the game is known as 𝑛 person game.
NULL HYPOTHESIS: The null hypothesis is the hypothesis that is being tested in a
hypothesis – testing situation. Often the null hypothesis is of the form. “There is no
relation between two quantities”.

NULL MATRIX OR ZERO MATRIX: The m × n matrix whose elements are all 0 is called
the null matrix (or zero matrix) of the type m × n. It is usually denoted by 𝑂 or more
clearly by 𝑂𝑚 ,𝑛 . Often a null matrix is simple denoted by the symbol 0 read as ‘zero’.

For example

0 0 0 0 0 0 0 0
0 0 0 0 0 and 0 0 0
0 0 0 0 0 3×5 0 0 0 3×3

Are zero matrices of type 3 × 5 and 3 × 3 respectively.

NULL SET: The null set is that contains no elements. The term “null set” mean the same
as the term “empty set”.

NULL TREE: A null tree is simply a tree with zero nodes.

NUMBER LINE: A number line is a line on which each point represents a real number.

NUMBER SYSTEM: Number system is most commonly used in Linear Algebra:

1. The Natural numbers ℕ = 1,2,3,4, ⋯ ⋯ . In ℕ, addition is closed but not subtraction;

e.g 2 − 3 ∈ ℕ

2. The integers ℞ = ⋯ ⋯ , −2, −1,0,1,2,3, ⋯ ⋯ . In 𝕫, addition, subtraction and

2
multiplication are always closed, but not division; e.g.3 ∈ 𝕫.

𝑝
3. The rational numbers ℚ = : 𝑝, 𝑞 𝜖 𝕫 , 𝑞 ≠ 0 . In ℚ, addition, subtraction,
𝑞

multiplication and division (except by zero) are all closed. However, 2 ∈ ℚ.

4. The Real numbers ℝ: These are the numbers which can be expressed as decimals. The
rational numbers are those with finite terminating or recurring decimals. In ℝ, addition,
subtraction, multiplication and division (except by zero) are still closed and all positive
numbers have square roots ,but −1 ∈ ℝ.
5. The complex numbers ℂ = 𝑥 + 𝑖𝑦: 𝑥, 𝑦 𝜖 ℝ , 𝑤𝑕𝑒𝑟𝑒 𝑖 2 = 1. In ℂ, addition, subtraction,
multiplication and division (except by zero) are still closed and all numbers have square
roots. In fact, all polynomial equations with coefficients in ℂ have solutions in ℂ.

NUMBER THEORY: Number theory is the study of properties of the natural numbers.
One aspect of number theory focuses on prime numbers. There are an infinite number
of prime numbers. Suppose for example that 𝑝 was the largest prime number. Then ,
form a new number equal to one plus the product of all the prime numbers from 2 up to
𝑝. This number will not be divisible by any of these prime numbers (and , therefore , not
by any composite number formed by multiplying these primes together) and will
therefore be prime. This contradicts the assumption that 𝑝 is the largest prime numbers.
There are still unsolved problems involving the frequency of occurrence of prime
numbers.

NUMERAL: A numeral is a symbol that stands for number For example “4” is the Arabic
numeral for the number four. “IV” is the Roman numeral for the same number.

NUMERATOR: The numerator is the number above the bar in a fraction . In the fraction
8
, 8 is the numerator.
9

NUMERICAL INTEGRATION: The numeral integration method is used when it is not

possible to find a formula that can be evaluated to give the value of a definite integral.

NUMERICAL METHOD OR ITERATIVE METHOD: Numerical methods concerns with the

iterative or trial and error methods. Whenever the classical methods fail, we use
iterative procedure. The classical methods may fail because of the complexity of the
constraints or of the number of variables.

In this procedure, we start with a trial solution and a set of rules for improving it. The
trial solution is improved by the given rules and is then replaced by this improved
solution. This process of improvement is repeated until no further improvement is
possible or when the cost of further calculation cannot be justified. Iterative procedures
as divided into three groups.

i) After a finite number of repetitions, no further improvement will be possible.

 ii) Although successive iterations improve the solutions, we are only guaranteed
the solution as a limit of an infinite process.
iii) Finally, we include the trial and error method which however is likely to be
lengthy, tedious and costly even if electronic computers are used.

𝒏-VECTOR SPACE: The set of all 𝑛-vectors of a field 𝐹 is called the 𝑛-vector space over 𝐹.
It is usually denoted by 𝑉𝑛 (𝐹) or simply by 𝑉𝑛 if the field is understood. Similarly the set
of all 3 vectors is a vector space which is usually denoted by 𝑉3 . The elements of the field
𝐹 are known as scalars relatively to the vectors.
O

OBJECTIVE FUNCTION: An objective function is a function whose value you are trying to
maximize or minimize. The value of the objective function depends on the values of a set
of choice variables and the problem is to find the optimal values for those choice
variables.

OBLATE SPHEROID: An oblate spheroid is elongated horizontally.

OCTAGON: Nan octagon is an eight-sided polygon.

OCTAHEDRON: An octahedron is a polyhedron with eight faces.

OCTAL: An octal number system is a base-eight number system.

ODD FUNCTION: The function 𝑓 𝑥 is an odd function if it is satisfies the property that
𝑓 −𝑥 = −𝑓 𝑥 . For example 𝑓 𝑥 = sin 𝑥 𝑎𝑛𝑑 𝑓 𝑥 = 𝑥 3 are both odd functions.

ODD NUMBER: An odd number is a whole number that is not divisible by 2, such as
1,3,5,7,9,11,13,15…

ODE: An ODE is an equation of the form 𝐺(𝑥, 𝑦, 𝑦 ′ , 𝑦 ′′ , . . . , 𝑦 (𝑛) ) = 0 regarded as an

equation for 𝑦(𝑥). We refer to 𝑦 as the dependent variable and 𝑥 as the independent
varaible. Usually this can be solved for the highest derivative of 𝑦 and written in the
form
𝑑𝑛 𝑦
𝑛
= 𝑦 (𝑛) (𝑥) = 𝐹(𝑥, 𝑦, 𝑦 ′ , 𝑦 ′′ , . . . , 𝑦 (𝑛−1) ).
𝑑𝑥
Then the order of the ODE is 𝑛, the order of the highest derivative which appears.
ONE-DIMENSIONAL, TWO-DIMENSIONAL AND THREE-DIMENSIONAL FLOWS: One-
dimensional flow neglects variations or changes in velocity, pressure etc, transverse to
the main flow direction. The flow characteristics vary only in the direction of flow.

The flow is said to the two-dimensional when the velocity vector q is everywhere at
right angles to a certain direction and independent of displacement parallel to that
direction. The components of q (𝑢, 𝑣, 𝑤) in a two-dimensional flow are (𝑢, 𝑣, 𝑤), where 𝑢
and 𝑣 are independent of 𝑧. The flow conditions are identical in planes perpendicular to
the 𝑧- axis.

In three-dimensional flow the velocity components (𝑢, 𝑣, 𝑤) in mutually perpendicular

direction are the functions of space coordinates and time 𝑥, 𝑦, 𝑧 and 𝑡.

ONE –TAILED TEST: Ina a one- tailed test the critical region consists of only one tail of a
distribution. The null hypothesis is rejected only if the test statistic has an extreme
value in one direction.

ONE–TO-ONE FUNCTION: A function 𝑦 = 𝑓 𝑥 is a one – to-one function if every value of

𝑥 in the domain is associated with a unique value of 𝑦 in the range, making it possible to
find an inverse function.

OPEN BALL: Let (𝑋, 𝑑) be a metric space and 𝑥0 ∇ 𝑋. Let 𝑟 be a positive number. The set
𝐵(𝑥0 , 𝑟) = {𝑥 ∇ 𝑋 ∶ 𝑑(𝑥, 𝑥0 ) < 𝑟} is called the ball with center 𝑥0 and radius 𝑟. On
some spaces like 𝐶 or 𝑅 2 this is also known as an open disk and when the space is 𝑅, it is
known as open interval.
OPEN INTERVAL: An open interval is an interval that does not contain both its
endpoints. For example, the interval 0 < 𝑥 < 1 is an open interval because the
endpoints 0 and 1 are not included.

OPEN MAPPING THEOREM (COMPLEX ANALYSIS): if 𝑈 is a domain of the complex

plane 𝑪 and 𝑓 : 𝑈 → 𝑪 is a non-constant holomorphic function, then 𝑓 is an open map.

OPEN MAPPING THEOREM (FUNCTIONAL ANALYSIS): If 𝑇 is a continuous

transformation of a Branch space 𝐵 onto a Banach space 𝐵 ′ , then 𝑇 is an open mapping.

OPERAND: An operand is a number that is the subject of an operation. In the equation

5+3=8, 5 and 3 are the operands.

OPERATION: An operation, such as addition or multiplication, is the process of carrying

out a particular rule on a set of numbers. The four fundamental arithmetic operations
are addition, multiplication, division and subtraction.
OPERATOR: An operator (or transformation) 𝑇 is an assignment of a vector 𝑇(𝑓) (also
written as 𝑇𝑓) to every vector 𝑓 in a certain subset 𝐷𝑇 of a normed linear space 𝑁. The
domain of definition of 𝑇 is 𝐷𝑇 and its range of values is denoted by ∆ 𝑇 .

OPTIMAL CODE: Let A be an alphabet of size 𝑞 > 1 and fix 𝑛, 𝑑. We define

𝐴𝑞 (𝑛, 𝑑) = 𝑚𝑎𝑥{𝑀 | 𝑡𝑕𝑒𝑟𝑒 𝑒𝑥𝑖𝑠𝑡𝑠 𝑎𝑛 (𝑛, 𝑀, 𝑑) − 𝑐𝑜𝑑𝑒 𝑜𝑣𝑒𝑟 𝐴}

An (𝑛, 𝑀, 𝑑)-code for which 𝑀 = 𝐴𝑞 (𝑛, 𝑑), is called an optimal code.

OPTIMAL FEASIBLE SOLUTION OF THE TRANSPORTATION PROBLEM: A feasible
solution (𝑥𝑖,𝑗 ) of the Transportation Problem is said to be optimal if it minimizes cost
amongst all feasible solutions of the Transportation Problem.
OPTIMAL LINEAR CODE: Let 𝑞 > 1 be a prime power and fix 𝑛, 𝑑. We define
𝐵𝑞 (𝑛, 𝑑) = 𝑚𝑎𝑥{𝑞 𝑘 | 𝑡𝑕𝑒𝑟𝑒 𝑒𝑥𝑖𝑠𝑡𝑠 𝑎 𝑙𝑖𝑛𝑒𝑎𝑟 [𝑛, 𝑘, 𝑑] − 𝑐𝑜𝑑𝑒 𝑜𝑣𝑒𝑟 𝐹𝑞𝑛 }
A linear [𝑛, 𝑘, 𝑑]-code for which 𝑞 𝑘 = 𝐵𝑞 (𝑛, 𝑑) is called an optimal linear code.
OPTIMAL SOLUTION: The optimal solution is that solution which is also a feasible
solution and in addition to it, it must optimizes (maximizes/ minimizes) the objective
function of the LPP. Consider a general linear programming problem with n real
variables 𝑥1 , 𝑥2 , . . . , 𝑥𝑛 whose objective is to maximize or minimize some objective
function subject to appropriate constraints. A optimal solution of this linear
programming problem is specified by an n-dimensional vector 𝑥 that is a feasible
solution that optimizes the value of the objective function amongst all feasible solutions
to the linear programming problem.

OPTIMAL STRATEGY AND VALUE OF THE GAME: If in a given payoff matrix 𝑖, 𝑗 is the
saddle point, then the player A and B are said to have the optimal strategy 𝑖 and 𝑗 . The
𝑡𝑕
value of the 𝑖, 𝑗 cell is known as the value of the game and is denoted by 𝑣.

OR: The word “OR” is a connective word used in logic. The sentence “𝑝 𝑜𝑟 𝑞" is false only
if both 𝑝 𝑎𝑛𝑑 𝑞 are false; it is true if either 𝑝 𝑎𝑛𝑑 𝑞 or both are true. The operation of OR
is illustrated by the truth table:

p q p or q
T T T
T F T
F T T
F F F
The symbol ∀ is often use to represent OR. An OR sentence is also called a disjunction.

ORBIT OF AN ELEMENT: Let 𝐿 be a field, let 𝐺 be a group of automorphisms of 𝐿, and let

𝛼 be an element of 𝐿. The orbit of 𝛼 under the action of 𝐺 on 𝐿 is the set {𝜍(𝛼) ∶ 𝜍 ∇ 𝐺}.
ORBIT-STABILIZER THEOREM: Given a group action 𝐺 on a set 𝑋, define 𝐺𝑥 to be the
orbit of 𝑥 and 𝐺𝑥 to be the set of stabilizers of 𝑥. For each 𝑥 ∇ 𝑋 the correspondence
𝑔(𝑥) → 𝑔𝐺𝑥 is a bijection between 𝐺𝑥 , and the set of left cosets of 𝐺𝑥 .
ORDERED PAIR: An ordered pair is a set of two numbers where the order in which the
numbers are written has an agreed –upon meaning. One common example of an
ordered pair is the Cartesian coordinates 𝑥, 𝑦 where it is agreed that the horizontal
coordinates is always listed first and the vertical coordinate last.

ORDERING RELATION: Let 𝑆 be a set. An ordering relation is a relation ≤ on 𝑆 such that,

for every 𝑎, 𝑏, 𝑐 ∇ 𝑆:
• Either 𝑎 ≤ 𝑏, or 𝑏 ≤ 𝑎,
• If 𝑎 ≤ 𝑏 and 𝑏 ≤ 𝑐, then 𝑎 ≤ 𝑐,
• If 𝑎 ≤ 𝑏 and 𝑏 ≤ 𝑎, then 𝑎 = 𝑏.
Given an ordering relation ≤, one can define a relation < by: 𝑎 < 𝑏 if 𝑎 ≤ 𝑏 and 𝑎 ≠
𝑏. The opposite ordering is the relation > given by: 𝑎 > 𝑏 if 𝑏 ≥ 𝑎, and the relation >
is defined analogously.
ORDER OF AN ELLIPTIC FUNCTION: The order of an elliptic function is defined as the
number of its poles in a cell.

ORDINATE: The ordinate of a point is another name for the 𝑦 coordinate.

ORIGIN: The origin is the point 0,0 in Cartesian coordinates. It is the point where the
𝑥 − 𝑎𝑥𝑖𝑠 and the 𝑦 − 𝑎𝑥𝑖𝑠 intersect.

ORTHOCENTER: The orthocenter of a triangle is the point where the three altitude of
the triangle meet.

ORTHOGONAL: Orthogonal means perpendicular.

ORTHOGONALITY OF JACOBI POLYNOMIALS: The Jacobi polynomials

Gn (x) are orthogonal w. r. t. weight function 𝑝 𝑥 = x q−1 1 − x p−q
. (q > 0, 𝑝 − 𝑞 >
−1)

On the integral 0 ≤ x ≤ 1

x1−q 1 − x q−p dn q+n−1 p+n−q

𝐺𝑛 𝑥 = [x 1−x ]
q q + 1 … (q + n − !) dx n

1
Since ∫0 𝑝 𝑥 . 𝐺𝑚 𝑥 𝐺𝑛 𝑥 𝑑𝑥 = 0.

ORTHOGONAL PROPERTIES OF HERMITE POLYNOMIALS: Orthogonal Properties of

Hermite Polynomials are given as

∞ 0
𝑥2 𝑖𝑓 𝑚 ≠ 𝑛
𝑒 − 𝐻𝑛 𝑥 𝐻𝑚 𝑥 𝑑𝑥 =
−∞ 𝜋22 𝑛 !
𝑖𝑓 𝑚 = 𝑛

ORTHOGONAL PROPERTIES OF LAGUERRE POLYNOMIALS: Orthogonal Properties of

Laguerre Polynomials are given as

∞
𝑥2 0 𝑖𝑓 𝑚 ≠ 𝑛
𝑒 − 𝐿𝑛 𝑥 𝐿𝑚 𝑥 𝑑𝑥 = 𝛿𝑚𝑛 =
−∞
1 𝑖𝑓 𝑚 = 𝑛

ORTHOGONAL PROPERTIES OF LEGENDRE’S POLYNOMIALS: Orthogonal Properties of

Legendre’s Polynomials are given as

+1
(i) ∫−1 Pm x Pn x dx = 0 if m ≠ n.
+1 2 2
(ii) ∫−1 Pn x dx = 2n+1 if m = n.

ORTHONORMAL: A set of vector is orthonormal if they are all orthogonal

(perpendicular) to each other, and they all have length l.

ORDER OR AN INTEGRAL FUNCTION (COMPLEX ANALYSIS): An integral function 𝑓(𝑧)

is said to be of finite order if there exists a positive number 𝑘, independent of 𝑟, such
that its maximum modulus 𝑀(𝑟) on the circle 𝑧 = 𝑟 satisfies the inequality.

log 𝑀 𝑟 < 𝑟 𝑘 for all sufficiently large values of 𝑟.

If there exists no such number 𝑘, then the integral function 𝑓(𝑧) is said to be of infinite
order.
ORIENTATION (MANIFOLD): Let V be a finite dimensional vector space. Two ordered
bases (𝑣1 , . . . , 𝑣𝑛 ) and (𝑣1 , . . . , 𝑣𝑛 ) are said to be equally oriented if the transition matrix
𝑆, whose columns are the coordinates of the vectors 𝑣1 , . . . , 𝑣𝑛 with respect to the basis
(𝑣1 , . . . , 𝑣𝑛 ), has positive determinant. Being equally oriented is an equivalence relation
among bases, for which there are precisely two equivalence classes. The space 𝑉 is said
to be oriented if a specific class has been chosen, this class is then called the orientation
of 𝑉, and its member bases are called positive. The Euclidean spaces 𝑅 𝑛 are usually
oriented by the class containing the standard basis (𝑒1 , . . . , 𝑒𝑛 ). For the null space
𝑉 = {0} we introduce the convention that an orientation is a choice between the signs
+ and −. An orientation of an abstract manifold 𝑀 is an orientation of each tangent
space 𝑇𝑝 𝑀, 𝑝 ∇ 𝑀, such that there exists an atlas of 𝑀 in which all charts induce the
given orientation on each tangent space. The manifold is called orientable if there exists
an orientation. If an orientation has been chosen we say that 𝑀 is an oriented manifold
and we call a chart positive if it induces the proper orientation on each tangent space.
A diffeomorphism map 𝑓: 𝑀 → 𝑁 between oriented manifolds of equal dimension is
said to be orientation preserving if for each 𝑝 ∇ 𝑀, the differential 𝑑𝑓𝑝 maps positive
bases for 𝑇𝑝 𝑀 to positive bases for 𝑇𝑓(𝑝) 𝑁. Every manifold 𝑀, of which there exists an
atlas with only one chart, is orientable. The orientation induced by that chart is of
course an orientation of 𝑀.
ORTHOGONAL COMPLEMENT: Let 𝑀 be a subspace of an inner product space V. The
orthogonal complement, written 𝑀⊥ , of 𝑀 is 𝑀⊥ = {𝑥 ∇ 𝑉: ⌌𝑥, 𝑚⌍ = 0 ∀𝑚 ∇ 𝑀.

 If 𝑀 is a closed subspace of a Hilbert space 𝐻 then 𝑀⊥ is a closed subspace

too (hence a Hilbert space too).
 Let 𝑀 be a closed subspace of a Hilber space 𝐻. Then for any 𝑥 ∇ 𝐻 there
exists the unique decomposition 𝑥 = 𝑚 + 𝑛 with 𝑚 ∇ 𝑀, 𝑛 ∇ 𝑀⊥ and
⊥
||𝑥||2 = ||𝑚||2 + ||𝑛||2 . Thus 𝐻 = 𝑀 ⊕ 𝑀 ⊥ and (𝑀⊥ ) =𝑀.
 The space 𝐶𝑃 −𝜋, 𝜋 is dense in 𝐿2 −𝜋, 𝜋 .
 It is not always true that ||𝑓𝑛 − 𝑓||∞ → 0 even for 𝑓 ∇ 𝐶𝑃[−𝜋, 𝜋].

ORTHOGONALLY SIMILAR MATRICES: Let 𝐴 and 𝐵 square matrices of order 𝑛. Then 𝐵 is

said to be orthogonally similar to 𝐴 if there exists an orthogonal matrix 𝑃 such that
𝐵 = 𝑃−1 𝐴𝑃. If 𝐴 and 𝐵 are orthogonally similar, then they are si
milar also. Further it can be easily shown that the relation of being ‘orthogonally similar’
is an equivalence relation in the set of all 𝑛 × 𝑛 matrices over the field of complex
numbers.

 Every real symmetric matrix is orthogonally similar to a diagonal matrix with

real elements.
 A real symmetric matrix of order 𝑛 has 𝑛 mutually orthogonal real eigenvectors.
 Any two eigenvectors corresponding to two distinct eigenvalues a real
symmetric matrix are orthogonal.
 If λ occurs exactly 𝑝 times as an eigenvalues of a real symmetric matrix 𝐴, then 𝐴
has 𝑝 but more then 𝑝 mutually orthogonal real eigenvectors corresponding to λ.

ORTHOGONAL MATRIX: A square matrix 𝐴 is said to be orthogonal if 𝐴′ 𝐴 = 𝐼.

If 𝐴, 𝐵 be 𝑛- rowed orthogonal matrices, AB and BA are also orthogonal matrices.

ORTHOGONAL PROPERTIES OF CHEBYSHEV POLYNOMIALS:

0 𝑚≠𝑛
1 Tm x Tn x
(i) ∫−1 [ 1−x 2 𝑑𝑥 = 𝜋/2 𝑚=𝑛≠0
0 𝑚=𝑛=0
0 𝑚≠𝑛
1 U m x U n (x)
(ii) ∫−1 𝑑𝑥 = 𝜋/2 𝑚=𝑛≠0
[ 1−x 2
0 𝑚=𝑛=0

ORTHOGONAL SYSTEM OF CURVES: Two families of curves u x, y = c1 , v x, y = c2 are

said to form an orthogonal system if they intersect at right angles at each of their points
of intersection.

ORTHOGONAL TRAJECTORIES: Consider the family of curves 𝜙 𝑢, 𝑣 = 𝑐

Lying on the surface 𝑟 = 𝑟 𝑢, 𝑣 . In case there exists another family of curves

Ψ 𝑢, 𝑣 = 𝑐1

Lying on the same surface such that at every point of the surface the two curves, one
from each family are orthogonal, then the family of curves (2) is called the orthogonal
trajectories of the family of curves (1)

ORTHOGONAL VECTORS: Two vectors 𝑥 and 𝑦 in an inner product space are said to be
orthogonal if ⟨ 𝑥, 𝑦 ⟩ = 0, written 𝑥 ⊥ 𝑦.
An orthogonal sequence (or orthogonal system) 𝑒𝑛 (finite or infinite) is one in which
𝑒𝑛 ⊥ 𝑒𝑚 whenever 𝑛 ≠ 𝑚.
An orthonormal sequence (or orthonormal system) 𝑒𝑛 is an orthogonal sequence with
||𝑒𝑛 || = 1 for all 𝑛.

If 𝑥 ⊥ 𝑥 then 𝑥 = 0 and consequently 𝑥 ⊥ 𝑦 for any 𝑦 ∇ 𝐻.

1. If all vectors of an orthogonal system are non-zero then they are linearly
independent.

2. Let 𝐴 be a subset of an inner product space 𝑉 and 𝑥 ⊥ 𝑦 for any 𝑦 ∇ 𝐴. Then

𝑥 ⊥ 𝑧 for all 𝑧 ∇ 𝐶𝐿(𝐴).

ORTHONORMAL TRIAD OF FUNDAMENTAL UNIT VECTORS 𝒕, 𝒏, 𝒃: We have defined a set

of three mutually perpendicular unit vectors 𝑡, 𝑛, 𝑏 associated with each point of a
curve. This set of unit orthonormal triad forms a moving trihedral at point 𝑃 (say) such
that
𝑡 ∙ 𝑛 = 0, 𝑛 ∙ 𝑏 = 0, 𝑏 ∙ 𝑡 = 0
𝑛 × 𝑏 = 𝑡, 𝑏 × 𝑡 = 𝑛, 𝑡 × 𝑛 = 𝑏.
The vectors 𝑡, 𝑛, 𝑏 are called fundamental unit vectors.
ORTHONORMAL VECTOR: A set 𝑀 of vectors of an inner product space 𝐻 is said to be
orthonormal if

(i) Each vector in 𝑀 is normalized, i.e. is a unit vector, and

(ii) Any two different vectors in 𝑀 are mutually orthogonal

We also say that the vectors of 𝑀 are orthonormal. If 𝑓, 𝑔 are two vectors of an
orthonormal set 𝑀, then 𝑓, 𝑔 = 0 for 𝑓 ≠ 𝑔 and 𝑓, 𝑔 = 1 𝑓𝑜𝑟 𝑓 = 𝑔

If 𝑀 consists of vectors 𝑒1 , where 𝑖 runs through a certain index set (which may be
finite, countably infinite, or even uncountable), then

𝑒𝑖 , 𝑒𝑗 = 0 for 𝑖 ≠ 𝑗, and 𝑒𝑖 , 𝑒𝑗 = 1 for 𝑖 = 𝑗

We also write this as 𝑒𝑖 , 𝑒𝑗 = 𝛿𝑖𝑗 , where 𝛿𝑖𝑗 known as the Kronecker delta, is
defined by

𝛿𝑖𝑗 = 0, for 𝑖 ≠ 𝑗,
 = 1, for 𝑖 = 𝑗.

A vector 𝑕 is said to be orthogonal to 𝑀, written as 𝑕 ⊥ 𝑀, if it is orthogonal to

every vector of 𝑀, 𝑖. 𝑒. , 𝑕, 𝑓 = 0 for all 𝑓 in 𝑀.

(iii) Equivalent generating system: Two arbitrary subsets 𝐴 and 𝐵 of a linear space
are said to be equivalent generating systems if the linear manifolds spanned
by them are equal, i.e., 𝐿𝑀 𝐴 = 𝐿𝑀 𝐵 . This means that any finite linear
combination of vector of 𝐴 can also be expressed as finite linear combination
of vectors of 𝐵, and conversely.

OSCULATING CIRCLE (OR THE CIRCLE OF CUREVATURE): Let 𝑃, 𝑄, 𝑅 be three points on

any curve then the circle of curvature at point 𝑃 is the limiting position of the circle
through 𝑃, 𝑄, 𝑅 which the points 𝑄, 𝑅 tend to 𝑃.
OSCULATING PLANE: The osculating plane at a point 𝑃 of a curve of class ≥ 2 is the
limiting position of the plane which contains the tangent line at 𝑃 and a neighboring
point 𝑄 on the curve as 𝑄 → 𝑃.
OSCULATING SPHERE (OR THE SPHERE OF CURVATURE): If 𝑃, 𝑄, 𝑅, 𝑆 are four points on
a curve then the sphere of curvature at point 𝑃 is the limiting position of the sphere
𝑃𝑄𝑅𝑆 when the points 𝑄, 𝑅, 𝑆 tend to coincide with 𝑃. Its radius and centre are called
radius and centre of spherical curvature.
OUTER MEASURE: Let 𝑆 be a semi-ring of subsets in 𝑋, and µ be a measure defined on 𝑆.
An outer measure µ∗ on X is a map µ∗ : 2𝑋 → [0, ∞] defined by:

µ∗ (𝐴) = 𝑘 𝜇 𝐴𝑘 ; 𝐴 ⊂ ⋃𝑘 𝐴𝑘 , 𝐴𝑘 ∇ 𝑆

An outer measure has the following properties:

1. µ∗ (∅) = 0;

2. 𝑖𝑓 𝐴 ⊆ 𝐵 𝑡𝑕𝑒𝑛 µ∗ (𝐴) ≤ µ∗ (𝐵);

3. if 𝐴𝑛 is any sequence in 2𝑋 , then µ∗ ∪ 𝑛 𝐴𝑛 ≤ 𝑛 µ∗ 𝐴𝑛 .

4. Let 𝑎 < 𝑏. Then µ∗ ([𝑎, 𝑏]) = 𝑏 − 𝑎.

PARABOLIC FRACTIONAL: A linear functional (bilinear) with one fixed points 𝑧0 is

called parabolic and is expressible as
 1 1
= + 𝑕 𝑖𝑓 𝑧0 ≠ ∞
𝑤 − 𝑧0 𝑧 − 𝑧0

Or 𝑤 = 𝑧 + 𝑕 𝑖𝑓 𝑧0 𝑖𝑓 𝑧0 = ∞

A linear fractional transformation with two different fixed points, 𝑧1 and 𝑧2 is

expressible as

𝑤 − 𝑧1 𝑘(𝑧 − 𝑧1 )
= 𝑖𝑓 𝑧1, 𝑧2 ≠ ∞
𝑤 − 𝑧2 (𝑧 − 𝑧2 )

If 𝑧2 = ∞, then it becomes 𝑤 − 𝑧1 = 𝑘(𝑧 − 𝑧1 )

A transformation with two different fixed points is called hyperbolic if 𝑘 > 0, and
elliptic if 𝑘 = 𝑎𝑒 𝑖𝛼 , where 𝑎 ≠ 1; 𝛼 and 𝑎 both are real number and 𝑎 > 0.

PARABOLIC POINTS (DIFFERENTIAL GEOMETRY): The points on the surface at which

the Gaussian curvature 𝐾 = 0 are called parabolic points. In this case 𝐿𝑁 − 𝑀² = 0

PARALLEL HYPERPLANE: Two hyperplane 𝐶1 𝑥 = 𝑍 and 𝐶2 𝑥 = 𝑍 are said to be parallel if

they have the same unit normals, i.e., if 𝐶1 = 𝜆𝑐2 for some 𝜆, 𝜆 being non – zero.

PARALLELOGRAM IDENTITY: In an inner product space H, we have

2 2 2 2
𝑥+𝑦 + 𝑥−𝑦 = 2{ 𝑥 + 𝑦 } for all 𝑥 and 𝑦 ∇ 𝐻

PARAMETRIZED MANIFOLD: A parametrized manifold in Rn is a smooth map

𝜍: 𝑈 → 𝑅 𝑛 , where 𝑈 ⊂ 𝑅 𝑚 is a non-empty open set. It is called regular at 𝑥 ∇ 𝑈 if the
𝑛 × 𝑚 Jacobi matrix 𝐷𝜍(𝑥) has rank m (that is, it has linearly independent columns),
and it is called regular if this is the case at all 𝑥 ∇ 𝑈. An m-dimensional parametrized
manifold is a parametrized manifold 𝜍: 𝑈 → 𝑅 𝑛 with 𝑈 ⊂ 𝑅 𝑚 , which is regular (that is,
regularity is implied at all points when we speak of the dimension). Clearly, a
parametrized manifold with 𝑚 = 2 and 𝑛 = 3 is the same as a parametrized surface,
and the notion of regularity is identical to the one introduced in ordinary Geometry. For
𝑚 = 1 there is a slight difference with the notion of parametrized curves, because in
ordinary Geometry, we have required a curve 𝛾: 𝐼 → 𝑅 𝑛 to be defined on an interval,
whereas here we are just assuming 𝑈 to be an open set in 𝑅. Of course there are open
sets in 𝑅, which are not intervals, for example the union of two disjoint open intervals.
Notice however, that if 𝛾: 𝑈 → 𝑅 𝑛 is a parametrized manifold with 𝑈 ⊂ 𝑅, then for
each 𝑡0 ∇ 𝑈 there exists an open interval 𝐼 around 𝑡0 in 𝑈, and the restriction of 𝛾 to
that interval is a parametrized curve in the previous sense.
PARETO CHART: A Pareto chart, also called a Pareto distribution diagram, is a
vertical bar graph in which values are plotted in decreasing order of relative frequency
from left to right. Pareto charts are extremely useful for analyzing what problems need
attention first because the taller bars on the chart, which represent frequency, clearly
illustrate which variables have the greatest cumulative effect on a given system.

The Pareto chart provides a graphic depiction of the Pareto principle, a theory
maintaining that 80% of the output in a given situation or system is produced by 20% of
the input.

The Pareto chart is one of the seven basic tools of quality control. The independent
variables on the chart are shown on the horizontal axis and the dependent variables are
portrayed as the heights of bars. A point-to-point graph, which shows the cumulative
relative frequency, may be superimposed on the bar graph. Because the values of the
statistical variables are placed in order of relative frequency, the graph clearly reveals
which factors have the greatest impact and where attention is likely to yield the greatest
benefit.

PARITY: The word parity applies to situations where two items or their properties may
be juxtaposed as being opposites (in a certain context) of each other. Integers are of
either odd or even parity when they are, respectively, odd or even. The convenience is
in being able to say "two numbers of different parities" without having to explicitly
mention which is which.
PARITY CHECK MATRIX (CODING THEORY): A parity check matrix 𝐻 for 𝐶 is a generator
matrix for the dual code 𝐶 ⊥ . A parity check matrix is said to be in standard form if it is of
the form (𝑌 | 𝐼𝑛−𝑘 ).
PARSEVAL’S FORMULA: If the functions 𝑓, 𝑔 ∇ 𝐿2 [−𝜋, 𝜋] have Fourier series
∞ 𝑖𝑛𝑡 ∞ 𝑖𝑛𝑡 𝜋 ∞
𝑓= 𝑛=−∞ 𝑐𝑛 𝑒 , 𝑔= 𝑛=−∞ 𝑑𝑛 𝑒 then ⌌𝑓, 𝑔⌍ ∫−𝜋 𝑓(𝑡)g(t) 𝑑𝑡 = 2𝜋 −∞ 𝑐𝑛 𝑑𝑛 . An
integrable function 𝑓 belongs to 𝐿2 [−𝜋, 𝜋] if and only if its Fourier series is convergent
and then ||𝑥||2 = 2𝜋 ∞
−∞ 𝑐𝑘 2 .
PARSEVEL’S IDENTITY FOR FOURIER TRANSFORM: If the complex fourier transform of
𝑓 𝑡 and 𝑔 𝑡 be 𝐹 𝑠 and 𝐺 𝑠 respectively, then

∞ ∞
(i) ∫−∞ 𝑓 𝑡 𝑔 𝑡 𝑑𝑡 = ∫−∞ 𝐹 𝑠 𝐺 𝑠 𝑑𝑠
Where 𝑔 𝑡 is complex conjugate of 𝑔 𝑡 𝑎𝑛𝑑 𝐺 𝑠 is complex conjugate of
𝐺 𝑠 .
∞ 𝟐 ∞ 𝟐
(ii) ∫−∞ 𝒇 𝑡 𝑑𝑡 = ∫−∞ 𝑭 𝑠 𝑑𝑠

PARSEVEL’S IDENTITY FOR FOURIER SINE AND COSINE TRANSFORM: If 𝐹𝑐 𝑠 and 𝐺𝑐 𝑠

are Fourier sine transforms and 𝐹𝑠 𝑠 𝑎𝑛𝑑 𝐺𝑠 𝑠 are Fourier cosine transforms of
𝑓 𝑡 𝑎𝑛𝑑 𝑔 𝑡 respectively,

Then

∞ 2 ∞
i. ∫𝟎 𝑓 𝑡 𝑔 𝑡 𝑑𝑡 = ∫𝟎 𝐹𝑠 𝑠 𝐺𝑠 𝒔 ds
𝜋
∞ 2 ∞
ii. ∫𝟎 𝑓 𝑡 𝑔 𝑡 𝑑𝑡 = ∫𝟎 𝐹𝐶 𝑠 𝐺𝐶 𝒔 ds
𝜋
∞ 2 2 ∞ 2
iii. ∫𝟎 𝑓 𝑡 𝑑𝑡 = ∫𝟎 𝐹𝑆 𝑠 ds
𝜋
∞ 2 2 ∞ 2
iv. ∫𝟎 𝑓 𝑡 𝑑𝑡 = ∫𝟎 𝐹𝐶 𝑠 ds
𝜋

PARSEVAL'S THEOREM: Suppose that 𝐴(𝑥) and 𝐵(𝑥) are two square integrable (with
respect to the Lebesgue measure), complex-valued functions on 𝑹 of period 2𝜋
with Fourier series

and

respectively. Then

where 𝑖 is the imaginary unit and horizontal bars indicate complex conjugation.
PARTIAL DIFFERENTIAL EQUATION: A Partial Differential Equation (PDE) is defined as
an equation relating a function of two or more independent variables and its partial
derivatives. The general Partial Differential equation is of the type
𝜕2𝑢 𝜕2𝑢 𝜕2𝑢 𝜕𝑢 𝜕𝑢
𝑎 𝜕𝑥 2 + 𝑏 𝜕𝑥𝜕𝑡 + 𝑐 𝜕𝑡 2 + 𝑑 𝜕𝑥 + 𝑒 𝜕𝑡 + 𝑓𝑢 = 𝑔 (1)

where 𝑎, 𝑏, 𝑐, 𝑑, 𝑒 𝑎𝑛𝑑 𝑓 may be functions of 𝑥, 𝑡 and even 𝑢.

Order of the Partial Differential Equation: The order of the equation is given by the
order of the highest derivative. Thus, if one of the functions 𝑎, 𝑏 or 𝑐 are non-zero, then
the PDE is of second order. If 𝑎 = 𝑏 = 𝑐 = 0but 𝑑 or 𝑒 are non-zero, then the PDE is of
first order. If the functions 𝑎, 𝑏, 𝑐, 𝑑, 𝑒 𝑎𝑛𝑑 𝑓 do not depend on the dependent variable 𝑢,
then equation (1) is called linear partial differential equation otherwise it is non-linear
partial differential equation. If 𝑔 = 0, then equation (1) is called homogeneous partial
differential equation otherwise it is non-homogeneous Partial Differential Equation.

Some examples of First order Partial Differential Equations:

𝜕𝑢 𝜕𝑢
1. − 𝜕𝑡 = 0 ( linear PDE)
𝜕𝑥
𝜕𝑢 𝜕𝑢
2. 𝑢 𝜕𝑥 − 𝜕𝑡 = 0 ( non-linear PDE)
𝜕𝑢 𝜕𝑢
3. 𝑒 𝑡 𝜕𝑥 − 2 𝜕𝑡 = 0 ( linear, non-homogeneous PDE)

Some examples of First order Partial Differential Equations:

𝜕2𝑢 𝜕𝑢
1. 2 𝜕𝑥 2 − 3 𝜕𝑡 = 0 ( linear PDE)
1 𝜕 𝜕𝑢 𝜕2𝑢
2. 𝑥 𝜕𝑥 − 𝜕𝑡 2 = 𝑥 2 (linear, non-homogeneous PDE)
𝑡 𝜕𝑥
𝜕2𝑢 𝜕2𝑢
3. − 𝑒𝑥 = 𝑢2 (non-linear PDE)
𝜕𝑥 2 𝜕𝑡 2

If 𝑢 = 𝑢(𝑟, 𝑠) and 𝑟 = 𝑟 𝑥, 𝑡 and 𝑠 = 𝑠 𝑥, 𝑡 , then the chain rule gives

𝜕𝑢 𝜕𝑢 𝜕𝑟 𝜕𝑢 𝜕𝑠 𝜕𝑢 𝜕𝑢 𝜕𝑟 𝜕𝑢 𝜕𝑠
= + 𝜕𝑠 and 𝜕𝑡 = + 𝜕𝑠
𝜕𝑥 𝜕𝑟 𝜕𝑥 𝜕𝑥 𝜕𝑟 𝜕𝑡 𝜕𝑡

It is to be noted that certain functions are solutions to first order PDEs. Consider the
first order wave equation
𝜕𝑢 𝜕𝑢
𝑐 − =0 (2)
𝜕𝑥 𝜕𝑡
where 𝑐 is a constant known as wave speed. Also, note that the following are all
solutions to Equation (2).
1. 𝑢 = 𝑠𝑖𝑛(𝑥 − 𝑐𝑡)
𝜕𝑢
2. = cos⁡
(𝑥 − 𝑐𝑡)
𝜕𝑥
𝜕𝑢
3. = −𝑐. cos⁡
(𝑥 − 𝑐𝑡)
𝜕𝑡

𝜕𝑢 𝜕𝑢
Consider the linear equation 𝑐 𝜕𝑥 + 𝜕𝑡 = −𝑢.

The solution is given by 𝑢 = 𝑒 −𝑡 𝐹(𝑥 − 𝑐𝑡), where 𝐹 is an arbitrary function.

𝜕𝑢 𝜕𝑢
Consider the linear equation without non-constant coefficients 𝑥 𝜕𝑥 + 𝜕𝑡 = 0, 𝑥 > 0.

The solution is given by 𝑢 = 𝐹(𝑥𝑒 −𝑡 ), where 𝐹 is again an arbitrary function but this
time 𝑧 = 𝑥𝑒 −𝑡 .

PARTIALLY ORDERED SET: A non-empty set on which some order relation is given.
Some examples of partially-ordered sets:
1) The set of natural numbers with the usual order relation.
2) The set of natural numbers, where 𝒂 ≤ 𝒃 means that 𝒂 divides 𝒃.
3) The set of all subsets of some set, where 𝒂 ≤ 𝒃 means that 𝒂 ⊆ 𝒃.
4) The set of all real-valued functions on the interval [𝟎, 𝟏], where 𝒇 ≤ 𝒈 means
that 𝒇(𝒕) ≤ 𝒈(𝒕) for all 𝒕 ∇ [𝟎, 𝟏].
5) The set of all finite increasing sequences of natural numbers, where
𝒂𝟏 , … , 𝒂𝒌 ≤ 𝒃𝟏 , … , 𝒃𝒍

means that 𝑘 ≤ 𝑙 and 𝑎𝑖 = 𝑏𝑖 for 1 ≤ 𝑖 ≤ 𝑘 .

6) An arbitrary non-empty set, where 𝒂 ≤ 𝒃 means that 𝒂 = 𝒃 (such a set is called a

trivial or discrete partially ordered set).

PARTITION: (a) Let 𝐼 = [𝑎, 𝑏] be a closed and bounded interval. Then by a partition of 𝐼,
we mean a finite set of real numbers 𝑃 = 𝑥0 , 𝑥1 , 𝑥2 , 𝑥3 , … , 𝑥𝑛 with the property that
𝑎 = 𝑥0 ≤ 𝑥1 ≤ 𝑥2 ≤ 𝑥3 ≤ ⋯ ≤ 𝑥𝑛 = 𝑏.
(b) Let A be a set. A partition of A is collection of subsets of A with the property that
every element of A belongs to exactly one of the subsets in the collection.
PARTITION OF NUMBERS: A Partition of a positive integer 𝑛 is an expression of 𝑛 as the
sum of positive integers. The numbers of partitions of 𝑛, where the order of the
summands is ignored and repetitions is permitted, is denoted by 𝑝(𝑛) and is called the
number of partitions of 𝑛. For example 𝑝 5 = 7 since 5 = 4 + 1 = 3 + 2 = 3 + 1 + 1 =
2 + 2 + 1 = 2 + 1 + 1 + 1 = 1 + 1 + 1 + 1 + 1. Therefore, 𝑝(𝑛) equal the number of
conjugate classes of the symmetric group of order 𝑛 and is closely related to the
representation theory of this group.

The generating function of 𝑝(𝑛) is

∞ ∞
𝑓 𝑥 =1+ 𝑛=1 𝑝 𝑛 𝑥𝑛 = 𝑛=1 1 − 𝑥𝑛 −1
.

PARTITION OF UNITY: Let 𝑀 be an abstract manifold. A partition of unity for

𝑀 is a collection (𝑓𝛼 )𝛼∇𝐴 of functions 𝑓𝛼 ∇ 𝐶 ∞ (𝑀) such that:
(1) 0 ≤ 𝑓𝛼 ≤ 1,
(2) the collection of the supports 𝑠𝑢𝑝𝑝 𝑓𝛼 , 𝛼 ∇ 𝐴, is locally finite in 𝑀,
(3) 𝛼 ∇𝐴 𝑓𝛼 (𝑥) = 1 for all 𝑥 ∇ 𝑀.
Notice that because of condition (2) the (possibly infinite) sum in (3) has only finitely
many non-zero terms for each 𝑥 (but the non-zero terms are not necessarily the same
for all 𝑥).
PATH-CONNECTED SETS: A subset 𝐸 of a metric space 𝑋 is called path-connected if the
following is true. To each pair of points 𝑎 and 𝑏 in 𝐸 corresponds a path 𝛾 in 𝐸 joining 𝑎
to 𝑏. Note that if 𝐸 is path-connected and 𝑓 is continuous on 𝐸 then 𝑓(𝐸) is path-
connected.

PATH LINE OF FLOW: The curve described in space by a moving fluid element is known
as its trajectory or path line. Such a lien is obtained by giving the position of an element
as a function of time. The path lines are given by 𝑞 = 𝑑𝑟/𝑑𝑡. In Cartesian coordinates.

𝑑𝑥 𝑑𝑦 𝑑𝑧
= 𝑢 𝑥, 𝑦, 𝑧, 𝑡 , = 𝑣 𝑥, 𝑦, 𝑧, 𝑡 , = 𝑤 𝑥, 𝑦, 𝑧, 𝑡 .
𝑑𝑡 𝑑𝑡 𝑑𝑡

In general, the path line varies with each fluid particle. It indicates the direction of
velocity of a single particle of fluid at different times. In the case of steady flow, the fluid
particle will travel along the stream line i.e., the path of a material element of fluid
coincides with a stream line.

PATHWISE CONNECTED: 𝑋 is called pathwise connected if for each pair of points

𝑎, 𝑏 ∇ 𝑆 there exists real numbers 𝛼 ≤ 𝛽 and a continuous map 𝛾: [𝛼, 𝛽] → 𝑋 such
that 𝛾(𝛼) = 𝑎 and 𝛾(𝛽) = 𝑏 (in which case we say that 𝑎 and 𝑏 can be joined by a
continuous path in 𝑋).
PAYOFF MATRIX: A’s payoff matrix

Player B

1 2 3 ⋯ 𝑛
𝑣11 𝑣12 𝑣13 … 𝑣1𝑛

𝑣21 𝑣22 𝑣23 … 𝑣2𝑛

𝑣31 ⋮ ⋮ ⋮

Player A ⋮ ⋮ ⋮ ⋮

B ’s payoff matrix 𝑣𝑚 1 𝑣𝑚 2 𝑣𝑚 3 … 𝑣𝑚𝑛

1 2 3 ⋯ 𝑛

− 𝑣11 − 𝑣12 − 𝑣13 … −𝑣1𝑛

−𝑣21 −𝑣22 − 𝑣23 … − 𝑣2𝑛

− 𝑣31 ⋮ ⋮ ⋮

⋮ ⋮ ⋮ ⋮

−𝑣𝑚 1 − 𝑣𝑚 2 − 𝑣𝑚 3 … − 𝑣𝑚𝑛

In 𝐴′ 𝑠 payoff matrix, the cell entry 𝑣𝑖𝑗 is the payment to player A when A choose the 𝑖 𝑡𝑕
activity and B choose the 𝑗 𝑡𝑕 .

Ina rectangular game or two person zero sum game, the player B’s payoff matrix will be
the negative of A’s payoff matrix. Thus the net gain will be zero.
PEANO EXISTENCE THEOREM: Peano existence theorem, Peano theorem or Cauchy-
Peano theorem is a fundamental theorem which guarantees the existence of solutions to
certain initial value problems.

Let D be an open subset of R × R with

a continuous function and

a continuous, explicit first-order differential equation defined on 𝐷, then every initial

value problem

for f with has a local solution

where is a neighbourhood of in , such that for all .

The solution need not be unique: one and the same initial value (x0,y0) may give rise to
many different solutions z.

PENDANT VERTEX: A vertex of a graph of degree 1 is said to be an pendant vertex.

PERFECT CODES: A code 𝐶 over an alphabet of size 𝑞 with parameters (𝑛, 𝑀, 𝑑) is called
a perfect code if
𝑞𝑛
𝑀=
𝑑−1
𝑉𝑞𝑛 2
Note that every perfect code is an optimal code, but not necessarily the other way
around.
PERIODIC FUNCTION: A function 𝑓 𝑧 is said to be a periodic function if there exists
dnon-zero constant 𝑇 such that.

𝐹 𝑧 = 𝐹 𝑧 + 𝑇 = 𝐹 𝑧 + 2𝑇 = 𝐹 𝑧 + 3𝑇 = − − −

For all values of 𝑧 this constant 𝑇 is called the period of the function 𝑓 𝑧 . Clearly if 𝑇 is a
period of 𝐹 𝑧 then 𝑛 𝑇 (where 𝑛 is a positive or negative integer) is also a period.

PERIOD OF A COMPLEX FUNCTION: A function 𝑓(𝑧) is said to be period 𝜔 if

𝑓 𝑧 + 𝜔 = 𝑓(𝑧)
The function ez is of period 2πi. For example, ez+2πi = 𝑒 𝑧 .

PERIPHERALLY CONTINUOUS: A function 𝑓 ∶ 𝑋 → 𝑌 is peripherally continuous,

𝑓 ∇ 𝑃𝐶(𝑋, 𝑌 ), if for every 𝑥 ∇ 𝑋 and for all pairs of open sets 𝑈 and 𝑉 containing 𝑥 and
𝑓 (𝑥), respectively, there exists an open subset 𝑊 of 𝑈 such that 𝑥 ∇ 𝑊 and
𝑓 [𝑏𝑑(𝑊 )] ⊂ 𝑉 , where 𝑏𝑑(𝑊) is the boundary of 𝑊. For the functions 𝑓 ∶ 𝑅 → 𝑅 this
means that 𝑓 has the Young property, that is, for every 𝑥 ∇ 𝑅 there exist sequences {𝑥𝑛 }
and {𝑦𝑛 } such that 𝑥𝑛 , 𝑦𝑛 , and both 𝑓 (𝑥𝑛 ) and 𝑓 (𝑦𝑛 ) converge to 𝑓 (𝑥). Young showed
that for the Baire class 1 functions, the Darboux property and the Young property are
equivalent.

PERMUTATIONS AND COMBINATIONS: Let there be given a set Ω and arrange them in a
row, we have a 𝑘-permutation of elements of Ω. The number of such arrays is
(𝑛)𝑘 = 𝑛 𝑛 − 1 … (𝑛 − 𝑘 + 1). The polynomial 𝑥 𝑘 = 𝑥 𝑥 − 1 … 𝑥 − 𝑘 + 1 in 𝑥 of
degree 𝑘 is called Jordan factorial of degree 𝑘. In particular, 𝑛 𝑘 = 𝑛!, n factorial, is the
number of permutations of Ω. A subset of Ω is called a 𝑘-subset if it contains exactly 𝑘
𝑛
elements. The number of 𝑘-subsets (or 𝑘-combinations) of Ω is 𝑘
= 𝑛 𝑘 /𝑘!. The
𝑥 𝑥
binomial coefficients 𝑛
are defined by the generating function 1 + 𝑧

∞ 𝑥
= 𝑛 =0 𝑛 𝑧 𝑛 . For any complex number 𝑥, the series is convergent for z < 1, and it is
𝑥
vertified that 𝑛
= 𝑥 𝑛 = 𝑛!in terms of the Jordan factorial 𝑥 𝑛 . The same result hold
in a comple field with valuation, in particular in a 𝑝-adic number field. in any case, we
have the recursive relation

𝑥 𝑥−1 𝑥−1
𝑛
= 𝑛
+ 𝑛−1
(𝑛 ≥ 1),

𝑥
0
= 1,

and in general

𝑛 𝑥 𝑦 𝑥+𝑦
𝑘=0 𝑘 = 𝑛−𝑘
= 𝑛
,

which leads to many identities involving binomial coefficients. The recursive relation
𝑛
allows us to compute the value of 𝑘
easily foe small integers 𝑛, 𝑘, as was noticed by
Pascal. The arrangement of these values in a triangular form:
 1
1 1
1 2 1
1 3 3 1
1 4 6 4 1

is called Pascal’s triangle. For integral values 𝑥, (1 + 𝑧)𝑥 are polynomials, and we have
𝑛 𝑛
(𝑎 + 𝑏)𝑛 = 𝑘=0 𝑘 𝑎𝑛 −𝑘 𝑏 𝑘 (binomial thorem). As a generalization, we have

𝑛
𝑛! 𝑝 𝑝
𝑎1 + ⋯ + 𝑎𝑚 = 𝑎1 1 … 𝑎𝑚𝑚
𝑝1 ! … 𝑝𝑚 !

(multinomial theorem), where the sum is extended over all nonnegative 𝑝𝑖 with
𝑝𝑖 = 𝑛.

The number of ways of choosing 𝑘 elements, allowing repetition from a set of 𝑛

−𝑛 𝑛+𝑘−1
elements is (−1)𝑘 𝑘
= 𝑘
. This is also the number of nonnegative integral
𝑛
solutions of 𝑖=1 𝑥𝑖 = 𝑘 . As an example of binomial coefficients with noninteger
−1/2 2𝑛
arguments, we have 𝑛
= −1 𝑛 2−2𝑛 .
𝑛

PERTURBATION: Let 𝑋 be a set and 𝑌 be a metric space. Let 𝐹, 𝐺 be two maps from 𝑋 to
𝑌 . We say that 𝐺 is a perturbation of 𝐹 if there exist a constant 𝑘 > 0 such that for each
𝑥0 , 𝑦0 ∇ 𝑋 one has:
𝑑(𝐺(𝑥0 ), 𝐺(𝑦0 )) ≤ 𝑘𝑑(𝐹(𝑥0 ), 𝐹(𝑦0 ))

We say that 𝐺 is a small perturbation of 𝐹 if it is a perturbation of constant 𝑘 < 1.

PHASE: Let the equation of the wave be taken as

y = a sin(mx − nt + ε)

where ε represents the phase of the wave at the instant from which t is measured. We
notice that wave motion have the same amplitude, wave length and period but they
differ in phase. The angle (𝑚𝑥 − 𝑛𝑡) is called the phase angle and n is called the phase
rate.

PHRAGMAN AND LINDELOFF THEOREM (COMPLEX ANALYSIS): Suppose 𝑓(𝑧) is

analytic and on a simple closed contour 𝐶 except at one point 𝑃 𝑜𝑓 𝐶. Let 𝑓(𝑧) ≤ 𝑀 on
𝐶 except at 𝑃. Suppose further that there is a function 𝑤 𝑧 , analytic and not zero in
𝐶 𝑠. 𝑡. 𝑤(𝑧) ≤ 1 inside 𝐶 and s.t., given any 𝜀 > 0, it is possible to obtain a system of
curves, arbitrary close to 𝑃 and connecting the two sides of 𝐶 𝑟𝑜𝑢𝑛𝑑 𝑃, on which
𝑤(𝑧) 𝑓(𝑧) ≤ 𝑀.

Pi 𝛑 : The ratio of the circumference of a circle to its diameter in a Euclidean plane is

denoted by π, the initial letter of περιμετροσ(perimeter). Thus 𝜋 can be defined as

1
2 ∫0 𝑑𝑥/ 1 − 𝑥 2 .

The symbol 𝜋 has been used since W.Jones (1675-1749) and L.Eular. The fact that this
ratio is a constant is stated in Euclid’s Elements; however, Euclid gave no statement
about the numerical value of 𝜋. As an approximate value of 𝜋, 3 has been used from
antiquity. According to the Rhind Papyrus, (4/3)4 was used in ancient Egypt. Let 𝐿𝑛 (𝑙𝑛 )
be the perimeter of a regular 𝑛-gen circumscribed about (inscribed in) a circle of radius
1. Then the relations

2 1 1
𝐿𝑛 > 𝜋 > 𝑙𝑛 , =𝐿 +𝑙 ,
𝐿2𝑛 𝑛 𝑛

𝑙2𝑛 = 𝑙𝑛 𝐿2𝑛

10 1
hold. Archimedes obtained 3 71 < 𝜋 < 3 7 by calculating 𝐿96 and 𝑙96 . In 3rd-century China

Liu Hui used 𝜋 = 3.14. In 5th –centry China, Tsu Chung-Chih mentioned 22/7 as in
inaccurate approximate value and 355/113 as an accurate approximate value of 𝜋.
These values were obtained by methods similar to those of Archimedes. In 5th-countury
India, Aryabhatta obtained 𝜋 = 3.1416, and in 16th country Europe, Adriaen van
Roomen obtained 𝜋 = 355/113.

F. Viete represented 2/𝜋 in the following infinite product:

∞
𝜋 1 1 1 1 1 1 1 1 1
cos = + × + + …
2n 2 2 2 2 2 2 2 2 2
n=2

Using this formula, L. van Ceulen (15440-1610) calculated 𝜋 to 35 decimals. In the 17th
and 18th centuries, the Japenese mathematicians 𝑇. Seki, K. Takebe, and Y. Matunage
computed 𝜋 to 50 decimals. Since the 17th century, many formulas that represent 𝜋 as a
sum of infinite series or as a limit have been used to obtain more accurate approximate
values. The following are representations of 𝜋 known in those days:

𝜋 2∙2∙4∙4∙6∙6…
= 1∙3∙3∙5∙5∙7∙7… (J. Wallis)
2

= 1 − 1/3 + 1/5 − 1/7 + ⋯ (J. Gregory, G.W.F. Leibniz)

= 4 Arctan 1/5- Arctan 1/239 (J. Machin).

A formula combining Machin’s representation of 𝜋 and the power series Arctan

𝑥 = 𝑥 − (1/3)𝑥 3 + 1 5 𝑥 5 − ⋯ is called Machin’s formula and was often used for
calculating an approximate value of 𝜋.

PICARD’S ‘GREAT’ THEOREM: (a) Let 𝐴 = {𝑧 ∇ 𝐶 ∶ 𝑅 < |𝑧| < ∞}, and let 𝑎, 𝑏, 𝑐 be
distinct elements of 𝐶 ∗ . Let 𝑓 ∶ 𝐴 → 𝐶 ∗ \{𝑎, 𝑏, 𝑐} be meromorphic. Then lim𝑧→∞ 𝑓(𝑧)
exists.

(b) Suppose that 𝑓 is transcendental and meromorphic in 𝐶. Then 𝑓 takes every value in
𝐶 ∗ , with at most 2 exceptions, infinitely often in 𝐶.

PICARD–LINDELÖF THEOREM: Picard–Lindelöf theorem, Picard's existence theorem or

Cauchy–Lipschitz theorem is an important theorem on existence and uniqueness of
solutions to first-order equations with given initial conditions.

Consider the initial value problem

Suppose  𝑓  is Lipschitz continuous in 𝑦 and continuous in 𝑡. Then, for some value 𝜀 >
0, there exists a unique solution 𝑦(𝑡)to the initial value problem on the interval 𝑡0 −
𝜖, 𝑡0 + 𝜖 .

PICARD’S ‘LITTLE’ THEOREM: Let 𝑓 ∶ 𝐶 → 𝐶\{0, 1} be analytic. Then 𝑓 is constant.

PICARD’S EXISTENCE THEOREM FOR NON-AUTONOMOUS SYSTEMS: Suppose that
𝑓 ∶ [−𝛿, 𝛿] × 𝑋 → 𝑋 is such that, for all 𝑥 ∇ [−𝛿, 𝛿], the field 𝑓(𝑥, 𝑢) is Lipschitz with
respect to 𝑢 on a closed ball, say 𝐵𝑅 (𝑢0 ),, of radius 𝑅 > 0 about
an element 𝑢0 ∇ 𝑋. Let
𝑀 = 𝑠𝑢𝑝 ||𝑓(𝑥, 𝑢)|| < ∞
𝑢∇𝐵𝑅 𝑢 0 ,𝑥 ∇[−𝛿,𝛿]

be an upper bound for 𝑓 . Then the initial value problem

 𝑢′ (𝑥) = 𝑓(𝑥, 𝑢(𝑥)), 𝑢(0) = 𝑢0
has a unique continuously differentiable local solution 𝑢(𝑥), defined in the time interval
−𝛿 < 𝑥 < 𝛿,, where
𝛿 = 𝑅/𝑀.
PICARD’S EXISTENCE THEOREM FOR AUTONOMOUS SYSTEMS: Suppose that the vector
function 𝑓 ∶ 𝑋 → 𝑋 is Lipschitz on a closed ball, say 𝐵𝑅 (𝑢0 ), of radius 𝑅 > 0 about an
element 𝑢0 ∇ 𝑋. Let
𝑀 = 𝑠𝑢𝑝 ||𝑓(𝑢)|| < ∞.
𝑢∇𝐵𝑅 (𝑢 0 )

Then the initial value problem

𝑢′ (𝑥) = 𝑓(𝑢(𝑥)), 𝑢(0) = 𝑢0
has a unique continuously differentiable local solution 𝑢(𝑥), defined in the interval
−𝛿 < 𝑥 < 𝛿, where
𝛿 = 𝑅/𝑀.
PICARD’S THEOREM: The ODE 𝑦 ′ = 𝑓(𝑥, 𝑦) with 𝑦(𝑎) = 𝑏 has a solution in the
rectangle 𝑅 ∶ |𝑥 − 𝑎| ≤ 𝑕, |𝑦 − 𝑏| ≤ 𝑘 provided:
 𝑓 is continuous in 𝑅, bounded by 𝑀 (so |𝑓(𝑥, 𝑦)| ≤ 𝑀) and 𝑀𝑕 ≤ 𝑘.
 𝑓 satisfies a Lipschitz condition in 𝑅.
Furthermore, this solution is unique.
PICARD’S THEOREM COMPLEX ANALYSIS): An integral function attains every finite
value with at most one possible exception.

PISANO: Leonardo Pisano, also known as Fibonacci was an Italian mathematician who
lived from 1170 - 1250. Fibonacci is sometimes called the greatest European
mathematician of the middle ages. As a boy, Fibonacci traveled, with his father, to North
Africa, where he learned the Arabic numeral system. Subsequently, Fibonacci studied
under leading Arabic mathematicians. Fibonacci is best known for introducing
mathematical concepts he learned in the Middle East to the Western world, including
the decimal system and the Fibonacci sequence , a series of numbers beginning with
zero or one and proceeding in such a manner that each number is the sum of the two
preceding numbers.

PLANE CURVE IN SPACE: A curve is known as plane curve if it lies on a plane. Otherwise
it is said to be a skew twisted or tortuous curve. The parametric equations of a plane
curve are 𝑥 = 𝑥(𝑡), 𝑦 = 𝑦(𝑡), 𝑧 = 𝑧(𝑡), where 𝑥, 𝑦, 𝑧 are real valued functions of a single
real parameter 𝑡 ranging over a set value 𝑎 ≤ 𝑡 ≤ 𝑏.

POINT-TO-POINT GRAPH: A point-to-point graph, also called a line graph, is a pictorial

rendition of data in which specific values of a function are plotted as dots on a
coordinate plane. Adjacent pairs of dots are connected by straight lines. In most point-
to-point graphs, the independent variable is rendered along the horizontal axis with
values increasing from left to right. The dependent variable is rendered along the
vertical axis with values increasing from the bottom up.

The illustration is a point-to-point graph showing the temperature in degrees Celsius as

a function of the local time in a fictitious town over the 24-hour period representing a
hypothetical midsummer day. Temperature readings, accurate to the nearest degree,
are taken at 0000 hours, 0300 hours, 0600 hours, 0900 hours and so on at 3-hour
intervals until the following midnight. The results are plotted as points. Finally, each
adjacent pair of points is connected by a solid, straight line segment.

Point-to-point graphs are commonly used for portraying functions in which time is the
independent variable. A point-to-point graph may be superimposed on a bar graph in a
specialized plot called a Pareto chart . Some point-to-point graphs allow plotting of
functions having multiple dependent variables, positive/negative variables and multi-
category variables.
POISSON ALGEBRA: An algebra, usually over the field of real or complex numbers,
equipped with a bilinear mapping satisfying the properties of the usual Poisson bracket
of functions. Let 𝐴 be an associative commutative algebra over a commutative ring, 𝑅. A
Poisson algebra structure on 𝐴 is defined by an Rbilinear, skew-symmetric mapping,
{ , } ∶ 𝐴 × 𝐴 → 𝐴 such that

(i) (𝐴, { , }) is a Lie algebra over 𝑅,

(ii) (ii) the Leibniz rule is satisfied, namely, {𝑎, 𝑏𝑐} = {𝑎, 𝑏}𝑐 +
𝑏{𝑎, 𝑐}, for all a, b, c ∇ A.

The element {𝑎, 𝑏} is called the Poisson bracket of 𝑎 and 𝑏. The main example is that of
the algebra of smooth functions on a Poisson manifold

POISSON LIMIT THEOREM: If , such that , then

POISSON RANDOM VARIABLE: X is a Poisson random variable with parameter 𝜆 if

𝑒 −𝜆 𝜆𝑥
𝑓𝑋 (𝑥) = , 𝑥 = {0, 1, 2, . . . }
𝑥!
POISSION’S INTEGRAL FORMULA COMPLEX ANALYSIS): If 𝑓(𝑧) is analytic within and
on a circle 𝐶 defined by 𝑧 = 𝑅 and if 𝑎 is any point 𝐶, then

1 𝑅 2 − 𝑎𝑎 𝑓 𝑧 𝑑𝑧
𝑓 𝑎 = ∫
2𝜋𝑖 𝐶 𝑧 − 𝑎 (𝑅2 − 𝑧𝑎)

Hence the Poisson’s formula

2𝜋
𝑖𝜃
1 𝑅 2 − 𝑟 2 𝑓 𝑅𝑒 𝑖ϕ 𝑑𝜙
𝑓 𝑟𝑒 =
2𝜋 0 𝑅² − 2𝑅𝑟 cos 𝜃 − 𝜙 + 𝑟²

Where 𝑎 = 𝑟𝑒 𝑖𝜃 is any point inside the circle 𝑧 = 𝑅.

POLAR COORDINATES: Polar coordinates provide a method of rendering graphs and

indicating the positions of points on a two-dimensional (2D) surface. The polar
coordinate system is employed in mathematics, physics, engineering, navigation,
robotics, and other sciences. The polar plane consists of a reference axis, or ray, that
emanates from a point called the origin. Positions or coordinates are determined
according to the distance or radius, from the origin, symbolized r , and the angle relative
to the reference axis, symbolized by the lowercase Greek theta ( 𝜃). In the most
common polar system, the reference ray points off toward the right, and angles are
measured counterclockwise from it. This scheme is preferred, and is used by
mathematicians, physicists, and engineers. In a less common scheme, the reference ray
points upward, and angles are measured clockwise from it. This method is sometimes
used by astronomers, navigators, military personnel, meteorologists, and robotics
engineers.

POLAR FORM OF CAUCHY- RIEMANN EQUATION: If 𝑓 𝑧 = 𝑢 + 𝑖𝑣 is an analytic function

∂u 1 ∂u ∂v
and 𝑧 = 𝑟𝑒 𝑖𝜃 where 𝑢, 𝑟, 𝜃 are all real the Cauchy-Riemann equations are = , =
∂r r ∂θ ∂r
1 ∂u
− r ∂θ

POLAR PLANE OF A POINT WITH RESPECT TO A GIVEN SHAPE: The polar plane of a
given point with respect to a sphere is the locus of points the tangent planes at which
pass through the given point.

POLISH SPACES, LUZIN SPACES, AND SUSLIN SPACES: A topological space

homeomorphic to a complete separable metric space is called a Polish space. A
subspace 𝐸 of a Polish space 𝑋 is Polish if and only if it is a 𝐺𝛿 -subset of 𝑋, i.e., a
countable intersection of open subsets of 𝑋 (Alexandrov and Urysohn). A Hausdorff
topological space 𝑋 is called a Luzin space (resp. Suslin space) if we can find a Polish
space 𝑆 and a continuous bijective (resp. surjective) mapping 𝑓: 𝑆 + 𝑋. Every Polish
space is a Luzin space and every Luzin space is a Suslin space.
POLYGONAL NUMBERS: Let 𝑚 be an integer greater than 3, and let 𝑎1 = 1, 𝑎𝑛+1 − 𝑎𝑛 =
𝑚 − 2 𝑛 + 1; (𝑛 = 1,2, 3. . . ). The sequence {𝑎𝑛 } forms the system of polygonal
1
numbers of order 𝑚. The general term of {𝑎𝑛 } is given by 𝑛 + 2 𝑚 − 2 𝑛2 − 𝑛 ; (𝑛 =

1,2, ). Such 𝑎𝑛 are said to be triangular numbers if 𝑚 = 3, square numbers if 𝑚 = 4,

and pentagonal numbers if 𝑚 = 5.
Polynomial interpolation: Polynomial interpolation is a method of estimating values
between known data points. When graphical data contains a gap, but data is available
on either side of the gap or at a few specific points within the gap, an estimate of values
within the gap can be made by interpolation.
The simplest method of interpolation is to draw straight lines between the known data
points and consider the function as the combination of those straight lines. This method,
called linear interpolation, usually introduces considerable error. A more precise
approach uses a polynomial function to connect the points.

POLYNOMIAL RINGS: Addition and multiplication in 𝑅 𝑋 are defined by

𝑎𝑖 𝑋 𝑖 + 𝑏𝑗 𝑋𝑗 = (𝑎𝑖 𝑏𝑖 )𝑋 𝑖 ,

( 𝑎𝑖 𝑋 𝑖 ) + ( 𝑏𝑗 𝑋𝑗 ) = 𝑘 𝑖+𝑙=𝑘 𝑎𝑖 𝑏𝑗 𝑋 𝑘 .

A polynomial 𝑓 𝑋 𝜖 𝑅[𝑋] can be regarded as a function of a commutative ring 𝑅’

containing 𝑅 into itself such that 𝑐 ↦ 𝑓(𝑐). In this sense, 𝑓 𝑋 + 𝑔(𝑋) and 𝑓 𝑋 𝑔(𝑋) are
the functions such that 𝑐 ↦ 𝑓 𝑐 + 𝑔(𝑐) and 𝑐 ↦ 𝑓 𝑐 𝑔(𝑐), respectively.

It hold that

deg(𝑓 𝑋 + 𝑔(𝑋) ≤ 𝑚𝑎𝑥 deg 𝑓 𝑋 , deg 𝑔 𝑋

deg 𝑓 𝑥 𝑔(𝑋) ≤ deg 𝑓 𝑋 + deg 𝑔 𝑋

If 𝑅 is an integral domain, then the letter inequality is an equality and therefore 𝑅[𝑋] is
an integral domain. For these inequalities and for convenience elsewhere, we define the
degree of 0 to be indefinite.

Assume that 𝑅 is a field. For given 𝑓, 𝑔𝜖𝑅 𝑋 deg 𝑔 ≥ 1 , we can find unique 𝑞, 𝑟𝜖𝑅[𝑋]
such that 𝑓 = 𝑔𝑞 + 𝑟 and deg 𝑟 < deg 𝑔 𝑜𝑟 𝑟 = 0 (division algorithum). This 𝑔 is called
the integral quotient of 𝑓 by 𝑔, and 𝑟 is called the remainder of 𝑓 divided by 𝑔. The same
fact remains true in the general 𝑅 𝑋 if 𝑔(𝑋) is monic.

POLYNOMIALS IN ONE VARIABLE: Let 𝑅 be a commutative ring and 𝑎0 , 𝑎1 , ⋯ , 𝑎𝑛

elements of 𝑅. An expression 𝑓(𝑥) of the form

𝑓 𝑥 = 𝑎0 + 𝑎1 𝑋 + ⋯ , 𝑎𝑛 𝑋 𝑛 (1)

is called a polynomial in a variable 𝑋 over 𝑅; if 𝑎𝑛 ≠ 0, the number 𝑛 is called the degree

of the polynomials 𝑓(𝑋) and is denoted by deg 𝑓 . If 𝑎𝑛 = 1, the polynomial (1) is called
a monic polynomial. The totality of polynomials in 𝑋 over 𝑅 forms a commutative e ring
with respect to ordinary addition and multiplication, it is called the ring of polynomials
( or the polynomial ring) of 𝑋 over 𝑅 and is denoted by 𝑅[𝑋]. We say that we adjoin 𝑋 to
𝑅 to obtain 𝑅[𝑋].

POLYNOMIALS IN SEVERAL VARIABLE: Let 𝑅[𝑋, 𝑌] denote the ring 𝑅 𝑋 [𝑌], namely, the
ring obtained by adjoining 𝑌 to 𝑅[𝑋].

An element of 𝑅[𝑋, 𝑌] can then be expressed as 𝑎𝜇𝑣 𝑋𝜇 𝑌 𝑣 . This expression is called a

polynomial in 𝑋 and 𝑌 over 𝑅. Generally, 𝑅 𝑋1 , ⋯ , 𝑋𝑚 = 𝑅 𝑋1 , ⋯ , 𝑋𝑚 −1 [𝑋𝑚 ] is called
the polynomials ring in 𝑚 variables ( on 𝑚 indeterminates) 𝑋1 , ⋯ , 𝑋𝑚 over 𝑅, and its
element

𝑣 𝑣 𝑣
𝑓 𝑋1 , 𝑋2 , ⋯ , 𝑋𝑚 = 𝑎𝑣1 𝑣2 ⋯𝑣𝑚 𝑋1 1 𝑋2 2 ⋯ , 𝑋𝑚𝑚 (1)

is called a polynomial in 𝑚 variable 𝑋1 , ⋯ , 𝑋𝑚 over 𝑅. We call each summand a term of

the polynomial 𝑎𝑣1 𝑣2 ⋯𝑣𝑚 the degree of this term. The greatest degree of terms is called
the degree of the polynomials 𝐹. The terms 𝑎0 ⋯ 0 of degree 0 is called the constant
terms of 𝐹. If a polynomial 𝐹 in 𝑋1 , ⋯ , 𝑋𝑚 is composed of terms of the same degree 𝑛,
then 𝐹 is called a homogeneous polynomial (or from) of degree 𝑛; a polynomials
𝑣 𝑣 𝑣
consisting of a single terms, such as 𝑎 𝑋1 1 𝑋2 2 ⋯ , 𝑋𝑚𝑚 is called a monomial.

PONTRYAGIN’S DUALITY THEOREM: If 𝐺 is a Hausdorff locally compact Abelian group,

then the map 𝛷 ∶ 𝐺 → 𝐺 defined by 𝛷(𝑥)(𝜒) = 𝜒(𝑥) for 𝑥 ∇ 𝐺 and 𝜒 ∇ 𝐺 is an

algebraic isomorphism and a topological homeomorphism.

POSITON VECTOR: If the vector 𝑂𝑃 represents the position of the point P in space
relative to the point 𝑶 , then 𝑂𝑃 is called the position vector of P referred to 𝑶 as
origin.

If we say that 𝑨 is the point 𝑟 , then we mean that the position vector of 𝑨 is 𝒓 with
respect to some given origin O.

POSITIVE CORRELATION: A positive correlation is a relationship between

two variables such that their values increase or decrease together. Correlation is
expressed on a range from +1 to -1, known as the correlation coefficent. In a perfect
positive correlation, expressed as +1, an increase or decrease in one variable always
predicts the same directional change for the second variable. If two variables sometimes
but not always change in tandem, the correlation is expressed as greater than zero but
less than +1. Values below zero express negative correlation: As the value of one
variable increases, the other decreases. Zero indicates a lack of correlation: There is no
tendency for the variables to fluctuate in tandem either positively or negatively.

POSITIVE FUNCTIONAL: A functional F is positive if for any non-negative function 𝑓 we

have 𝐹(𝑓) > 0. Any positive linear functional 𝐹 on 𝐶(𝑋) is continuous and ||𝐹|| = 𝐹(𝟏),
where 𝟏 is the function identically equal to 1 on 𝑋.

Let λ be a continuous linear functional on 𝐶(𝑋). Then there are positive functionals 𝜆+
and 𝜆− on 𝐶(𝑋), such that 𝜆 = 𝜆+ − 𝜆−

For 𝑓 ∇ 𝐶ℝ (𝐾), we define the support of 𝑓, written as supp(𝑓), to be the closure of the
set {𝑥 ∇ 𝐾 ∶ 𝑓(𝑥) ≠ 0}.

POSITIVE OPERATORS IN 𝑪(𝑲)): Let 𝐶(𝐾) be the set of real-valued continuous maps on
a compact 𝐾. For those, there is a natural (partial) order: 𝑓 ≥ 𝑔 means 𝑓(𝑥) ≥ 𝑔(𝑥)
for all 𝑥 ∇ 𝐾. An operator 𝑈 ∶ 𝐶(𝐾) → 𝐶(𝐾) is called positive if 𝑓 ≥ 0 implies
𝑈(𝑓) ≥ 0 and it is called monotone if 𝑓 ≥ 𝑔 implies 𝑈(𝑓) ≥ 𝑈(𝑔). If 𝑈 is linear then it
is positive iff it is monotone.
POSITIVE ORIENTATION: The positive orientation (counterclockwise) is the one that is
given by the standard parametrization 𝑧(𝑡) = 𝑧0 + 𝑟𝑒 𝑖𝑡 , where 𝑡 ∇ 0, 2𝜋 .
POSTULATIONS OF LINEAR PROGRAMMING:
 Proportionality postulation: The contribution of each movement to the value
of the objective function is proportional to the level of that activity. In the
same way, the contribution of each movement to the left hand side of each
functional constraint is proportional to the level of the movement.
 Divisibility postulation: Decision variables are allowed to have any values,
including non-integer values that satisfy the functional and non-negativity
constraints.
 Certainty postulation: The value assigned to each parameter is assumed to be
a known constant.
 Additivity postulation: Every function is the sum of the individual
contributions of the relevant activities.
PRANDTL’S FUNCTION: The Prandtl’s stress function 𝜓 is defined as
1 2
𝜓 = 𝜓 𝑥, 𝑦 − 𝑥 + 𝑦2
2

𝜕𝜓 𝜕𝜓 𝜕𝜓 𝜕𝜓
so that = 𝜕𝑥 − 𝑥 , 𝜕𝑦 = 𝜕𝑦 − 𝑦. From these equations, we get
𝜕𝑥

𝜕𝜙 𝜕𝜓 𝜕𝜙 𝜕𝜓
= 𝜕𝑦 − 𝜕𝑦 = − and 𝜏𝑥𝑥 = 𝜏𝑦𝑦 = 𝜏𝑧𝑧 = 𝜏𝑥𝑦 = 0,
𝜕𝑥 𝜕𝑥

𝜕𝜙 𝜕𝜙
𝜏𝑧𝑥 = 𝜇 𝛼 − 𝑦 , 𝜏𝑦𝑧 = 𝜇 𝛼 +𝑥 ,
𝜕𝑥 𝜕𝑦

giving

𝜕𝜓
𝜏𝑧𝑥 = 𝜇 𝛼 ,
𝜕𝑦
𝜕𝜓
𝜏𝑧𝑦 = −𝜇 𝛼
𝜕𝑥

Since the stress components 𝜏𝑧𝑥 and 𝜏𝑧𝑦 are obtained from the function 𝜓 𝑥, 𝑦 by
differentiation , 𝜓 𝑥, 𝑦 is called the stress- function.

PREDICATE LOGIC: Predicate logic is the area of symbolic logic in which we take
quantifiers in account. Mainly propositional functions are discussed in predicate logic.
In the strict sense only single variable propositional functions are called predicates, but
the phrase predicate of n arguments (or wary predicate) denoting an 𝑛 variable
propositional function is also employed. Single-variable (or unary) predicates are also
called properties. We say that 𝑢 has the property 𝐹 if the proposition 𝐹(𝑎) formed by
the property 𝐹 is true. Predicates of two arguments are called binary relations. The
proposition 𝑅(𝑎, 𝑏) formed by the binary relation 𝑅 is occasionally expressed in the
form 𝑎𝑅𝑏. Generally, predicates of 𝑛 arguments are called n-ary relations. The domain of
definition of a unary predicate is called the object domain, elements of the object
domain are called objects, and any variable running over the object domain is called an
object variable. We assume here that the object domain is not empty. When we deal
with a number of predicates simultaneously (with different numbers of variables), it is
usual to arrange things so that all the independent variables have the same object
domain by suitably extending their object domains.
PRE-IMAGE THEOREM: Let 𝑓 ∶ 𝑀 → 𝑁 be a smooth map and let 𝑄 ∇ 𝑁 such that 𝑓 −1 (𝑄)
is not empty. Suppose that 𝑄 is regular. Then 𝑓 −1 (𝑄) is a submanifold of 𝑀 with dim
𝑓 −1 (𝑄)= dim𝑀 − dim𝑁.
PRESERVANCE OF CROSS RATIO: Cross ratio remains invariant under 𝑎 bilinear
transformation.

PRESSURE IN THE FLUID: The pressure P at a point in the fluid is the limit of the ratio of
normal force δA by the surrounding fluid particles as the area approaches zero. It is
defined as

Lt δF
p= .
δA → 0 δA

We know that the pressure at every point of an ideal fluid is equal in all directions
whether the fluid be at rest or in motion. It follows that an elements δA of a very small
area, free to rotate about its centre will have a force of constant magnitude acting on
either side of it.

PRIMARY DECOMPOSITION THEOREM: If 𝑀 is a finitely generated module over a

Euclidean domain 𝐷 then 𝑀 may be written as an internal direct sum 𝑀 = 𝑀1 ⊕ 𝑀2 ⊕
· · · ⊕ 𝑀𝑠 where 𝑀𝑖 are primary modules or free cyclic modules.
PRIME ELEMENT: An element 𝑝 in a ring 𝑅 is a prime element if it generates a prime
ideal. If 𝑅 is commutative, this is equivalent to saying that for all 𝑎, 𝑏 ∇ 𝑅 , if p divides
𝑎𝑏, then 𝑝 divides 𝑎 or 𝑝 divides 𝑏.

PRIME IDEALS: Let 𝑅 be a unital ring. A proper ideal 𝐼 is said to be prime if, given any
ideals 𝐽 and 𝐾 satisfying 𝐽𝐾 ⊂ 𝐼, either 𝐽 ⊂ 𝐼 or 𝐾 ⊂ 𝐼.

PRIME NUMBER: A prime number is an integer 𝑝 greater than one with the property
that 1 and 𝑝 are the only positive integers that divide 𝑝.

PRIME NUMBER THEOREM: Prime number theorem (PNT) describes

the asymptotic distribution of the prime numbers. The prime number theorem gives a
general description of how the primes are distributed among the positive integers. It
formalizes the intuitive idea that primes become less common as they become larger.

Let 𝜋(𝑥) be the prime-counting function that gives the number of primes less than or
equal to 𝑥, for any real number 𝑥. For example, 𝜋(10) = 4 because there are four prime
numbers (2, 3, 5 and 7) less than or equal to 10. The prime number theorem then states
that 𝑥 / 𝑙𝑛(𝑥) is a good approximation to 𝜋(𝑥), in the sense that the limit of the
quotient of the two functions 𝜋(𝑥) and 𝑥 / 𝑙𝑛(𝑥) as x approaches infinity is 1:

known as the asymptotic law of distribution of prime numbers. Using asymptotic

notation this result can be restated as

This notation (and the theorem) does not say anything about the limit of
the difference of the two functions as x approaches infinity. Instead, the theorem states
that 𝑥/𝑙𝑛(𝑥) approximates 𝜋(𝑥) in the sense that the relative error of this
approximation approaches 0 as x approaches infinity.

The prime number theorem is equivalent to the statement that the 𝑛th prime
number 𝑝𝑛 satisfies

,
The prime number theorem is also equivalent to

,
and

.
PRIMITIVE ELEMENT THEOREM: Every finite separable field extension is simple.

PRINCIPAL AND GENERAL VALUES OF AN IMPROPER INTEGRAL: Consider the integral

b
𝑓 𝑥 dx
a

in which 𝑓 𝑥 , become infinite at 𝑥 = 𝑐, where 𝑐 lies between the limits 𝑎 and 𝑏 . In this
case we define the integral

b
𝑓 𝑥 dx
a
as
c−μϵ b
lim 𝑓 𝑥 dx + 𝑓 𝑥 dx
∇→0 a c+υϵ

𝜇 and 𝑣 being arbitrary constants.

The above limit gives what is called the general value of the improper integral. If 𝜇 = 𝑣,
the value of the above limits is called the principal value of the integral.

PRINCIPAL AXIS THEOREM: The principal axis theorem concern quadratic forms in 𝑹𝒏 ,
which are homogeneous polynomials of degree 2. Any quadratic form may be
represented as

where A is a symmetric matrix.

PRINCIPAL BUNDLE: A principal bundle is a fiber bundle 𝑃 → 𝐵 together with

an action on 𝑃 by a Lie group 𝐺 that preserves the fibers of 𝑃 and acts simply
transitively on those fibers.

PRINCIPAL FRACTIONAL IDEAL: A principal fractional ideal of an integral domain 𝑅 is

an ideal of the form 𝑐𝑅, where 𝑐 is an element of the field of fractions of 𝑅. If R is not a
field, then the only fractional ideal of 𝑅 that is an ideal of 𝐾 is the zero ideal. When 𝑅 is
itself a field, the only fractional ideals of 𝑅 are the zero ideal and the whole of 𝑅, and
these are ideals of 𝑅. A fractional ideal of an integral domain 𝑅 is an 𝑅-module, and is
isomorphic as an 𝑅-module to some ideal of 𝑅. Indeed, given a fractional ideal 𝑀 of 𝑅
contained in the field of fractions of 𝑅, there exists a non-zero element 𝑎 of 𝑅 for which
𝑎𝑀 ⊂ 𝑅. Then 𝑎𝑀 is an ideal of 𝑅, and multiplication by a provides an isomorphism of
𝑅-modules between the fractional ideal 𝑀 of 𝑅 and the ideal 𝑎𝑀 of 𝑅.

PRINCIPAL IDEAL DOMAINS: An integral domain 𝑅 is said to be a principal ideal domain

(or PID) if every ideal of 𝑅 is a principal ideal.

PRINCIPAL NORMAL: Let 𝑡 𝑖 be unit tangent vector at any point 𝑃(𝑥 𝑖 ) lying on a curve 𝐶
in a Vn. The derived vector of 𝑡 𝑖 in its own direction is called fist curvature vector of C
relative to Vn and is denoted by 𝑝𝑖 . Then
 dx i dx i
𝑝𝑖 = 𝑡 𝑖 , where 𝑡 𝑖 =
dt ds

dx j 𝜕𝑡 𝑖 i dx i 𝜕𝑡 𝑖 dx i dx i i dt i dx a dx i i
pi = ti,𝑗 = + tα = 𝜕𝑥 𝑗 + tα aj
= + aj
ds 𝜕𝑥 𝑖 aj ds ds ds ds ds ds

d2xi dx i dx k i
= ds 2 + jk
ds ds

This is the required expression for 𝑝𝑖 . The magnitude of 𝑝𝑖 , denoted by 𝑘, is defined as

first curvature ( or curvature simple) of the curve C relative to Vn. Then 𝑘 2 = 𝑔𝑖𝑗 𝑝𝑖 𝑝 𝑗 . If
𝑛 be unit vector along 𝑝𝑖 , then we have 𝑝𝑖 = kni, ni being contravariant component of 𝑛. 𝑛
is called unit principal normal (geodesic normal) at the point P. The direction of 𝑝𝑖 , i.e.,
the unit vector ni, is called principal normal direction at 𝑃.

PRINCIPAL PART OF A FUNCTION: When 𝑓 has a pole at 𝑃, it is customary to call the

negative power part of the Laurent expansion of 𝑓 around 𝑃 the principal part of 𝑓 at 𝑃.

PRINCIPAL STRAIN: Let us consider infinitesimal affine deformation characterized by

the direction ratios of the lines through 𝑃0 𝑥 0 whose direction is left unchanged by the
deformation 𝛿𝐴𝑖 = 𝑒𝑖𝑗 𝐴𝑗 .

If the direction of vector Α is not changed by the strain, then δΑ and Α are parallel and
their components are proportional.

i.e. 𝛿𝐴𝑖 ∝ 𝐴𝑖

or 𝛿𝐴𝑖 = ℯ𝐴𝑖 𝑖 = 1,2,3

where e is the constant of proportionality.

𝛿𝐴𝑖
Since ℯ = is the extension of each component of vector Α and is thus the extension of
𝐴𝑖

Α itself.

𝛿Α
Or ℯ = ,
Α

From equation ℯΑ2 = 𝑒𝑖𝑗 𝑥𝑖 𝑥𝑗 , the extension ℯ is given by

𝑒𝑖𝑗 𝑥𝑖 𝑥𝑗
ℯ= ,
Α2
Also since , 𝛿𝐴𝑖 = ℯ𝑖𝑗 𝐴𝑗 ,

Or ℯ𝑖𝑗 𝐴𝑗 = 𝛿𝐴𝑖 = ℯ𝐴𝑖 = ℯ𝛿𝑖𝑗 𝐴𝑗 …. (1)

= ℯ𝑖𝑗 𝐴𝑗 = ℯ𝛿𝑖𝑗 𝐴𝑗 = ℯ𝑖𝑗 − ℯ𝛿𝑖𝑗 𝐴𝑗 =0 …. 2

Equation (2) is the set of three linear homogenous equations in the unknown 𝐴1 , 𝐴2 ,
𝐴3 . In expanded form, the equations (1) are

ℯ11 − ℯ 𝐴1 + ℯ12 Α2 + +ℯ13 Α3 = 0,

ℯ21 Α1 + ℯ22 − ℯ 𝐴2 + ℯ23 Α3 = 0

and ℯ31 Α1 + ℯ32 Α2 + ℯ33 − ℯ 𝐴3 = 0 …. 3

(3) possesses a non- vanishing solution if and only if the determinant of the coefficients
of 𝐴1 𝐴2 ,and 𝐴3 is equal to zero

𝑒11 − 𝑒 𝑒12 𝑒13

i.e. 𝑒21 𝑒22 − 𝑒 𝑒23 or ℯ𝑖𝑗 − ℯ𝛿𝑖𝑗 = 0
𝑒31 𝑒32 𝑒33 − 𝑒

Since this is cubic in ℯ and has three roots say ℯ1 , ℯ2 𝑎𝑛𝑑 ℯ3 . These three roots are
termed as the principal strains.

PRINCIPAL STRUCTURE THEOREM: If 𝐺 is a Hausdorff locally compact Abelian group

then we can find an open (so closed) subgroup 𝐻 such that 𝐻 = 𝑊 ⊕ 𝑅 𝑛 with 𝑊 a
compact group.

PRINCIPLE OF DOMINANCE: If one pure strategy of a player is better or superior than

another one, then the inferior strategy may be simply ignored by assigning a zero
probability while searching for optimal strategies.

PRINCIPLE OF LEAST ACTION: If 𝑇 be the kinetic energy at time 𝑡 of a Conservative

holonomic dynamical system specified by the generalized coordinates , then the integral
𝑡2
𝐼= 2𝑇𝑑𝑡
𝑡1
Has necessarily an extreme value, minimum or maximum , on actual paths as compared
with varied path as the system passes from one configuration at time to another
configuration at time 𝑡1 .
Consider all the curves passing through two fixed points P0 and P1 . If any of these curves
for which the distance P0 P1 measured along the curve is stationary, then that curve is
called geodesic.

PRINCIPLE OF SUPERPOSITION: To get the general solution we just add together all the
∞
solutions we have (𝑥, 𝑡) = 𝑛=1 sinnx(𝑎𝑛 𝑐𝑜𝑠𝐾𝑛𝑡 + 𝑏𝑛 𝑠𝑖𝑛𝐾𝑛𝑡) . Applying the initial
∞
condition 𝑦(𝑥, 0) = 𝐹(𝑥), this gives 𝐹 𝑥 = 𝑛=1 𝑎𝑛 𝑠𝑖𝑛𝑛𝑥 .

The right-hand side is a series of odd functions so if we extend 𝐹 to a function 𝐺 by

reflection in the origin, we have

𝐹 𝑥 , 𝑖𝑓 0 ≤ 𝑥 ≤ 𝜋;
𝐺 𝑥 =
−𝐹 −𝑥 , 𝑖𝑓 − 𝜋 ≤ 𝑥 < 0

∞
we have 𝐺 𝑥 = 𝑛=1 𝑎𝑛 𝑠𝑖𝑛𝑛𝑥 𝑓𝑜𝑟 − 𝜋 ≤ 𝑥 < 𝜋.

If we multiply by 𝑠𝑖𝑛𝑟𝑥 and integrate term by term, we get

1 𝜋
𝑎𝑘 = 𝜋 ∫−𝜋 𝐺 𝑥 𝑠𝑖𝑛𝑘𝑥𝑑𝑥,

This implies that the 𝑎𝑛 are precisely the sine coefficients of 𝐺. If we now assume,
further, that the right-hand side is differentiable term by term, we differentiate with
respect to 𝑡, and we set 𝑡 = 0, to get

∞
0 = 𝑦𝑡 (x, 0) = 𝑛=1 𝑏𝑛 𝐾𝑛𝑠𝑖𝑛𝑛𝑥

This equation is solved by taking 𝑏𝑛 = 0 for all 𝑛, so we have the following result.

A solution of the differential equation (1) with the given boundary and initial conditions
∞
is 𝑦 𝑥, 𝑡 = 𝑛=1 𝑎𝑛 𝑠𝑖𝑛𝑛𝑥𝑐𝑜𝑠𝐾𝑛𝑡, where the coefficients an are the Fourier sine
1 𝜋
coefficients 𝑎𝑘 = 𝜋 ∫−𝜋 𝐺 𝑥 𝑠𝑖𝑛𝑘𝑥𝑑𝑥 of the 2π periodic function G, defined on (−𝜋, 𝜋]

by reflecting the graph of F in the origin.

PROBABILISTIC MODEL: The inventory models, in which the demand is a random
variable having a known probabilistic distribution, are known as probabilistic models.
In probabilistic model, the future demand is determined by collecting data from the past
experience.

PROBABILITY SPACES: Let Ω be an abstract space and 𝔅 be a 𝜍- algebra of subsets of Ω.

A probability measure (or probability distribution) over Ω(𝔅) is a set function 𝑃(𝐸)
defined for 𝐸𝜀𝔅 and satisfying the following conditions: 𝑃1 𝑃 𝐸 ≥ 0; (𝑃2) for every
sequence 𝐸𝑛 (𝑛 = 1,2, … ) of pairwise disjoint sets in 𝔅,

𝑃 𝐸𝑛 = 𝑃 𝐸𝑛 ;
𝑛 𝑛

𝑃3 𝑃 Ω = 1. The triple Ω, 𝔅, 𝑃 is called a probability space. The space Ω (resp. each

element 𝜔 of Ω) is called the basic space, space of element events, or sample space
(resp. sample point or elementary event). We say that a condition 𝜀 𝜔 involving a
generic sample point 𝜔 is an event; in particular, it is called a measurable event or
random event if the set 𝐸 of all sample points satisfying 𝜀 𝜔 belongs to 𝔅. We assume
that an event is always a measurable one, since we encounter only measurable events in
the theory of probability. Because of the obvious one-to one correspondence between
measurable events and 𝔅- measurable sets (i.e., the correspondence of each event 𝜀
with the set 𝐸 of all sample points 𝜔 satisfying 𝜀), a 𝔅-measurable set itself is frequently
called an event. If 𝜀 is an event and 𝐸 is the 𝔅-measurable set corresponding to 𝜀, we call
𝑃 𝐸 or Pr(𝜀) the probability that the event 𝜀 occures, i.e., the probability of the event 𝐸.
The complementary event (resp. impossible event, sure event) is the complementary
set 𝐸 𝑐 (empty set ⊘, whole space Ω). For a finite or infinite family 𝐸𝜆 𝜆𝜖Λ , the sum
event (resp. intersection or product event) or 𝐸𝜆 is the set ⋃𝜆 𝐸𝜆 ⋂𝜆 𝐸𝜆 . If 𝐸 ⋂ 𝐹 =⊘,
then we say that 𝐸 and 𝐹 are mutually exclusive or that they are exclusive events.

By the definition of 𝑃, we have 0 ≤ 𝑃(𝐸) ≤ 1 for any event 𝐸, 𝑃(⊘) = 0, and 𝑃 Ω = 1.

Moreover, if 𝐸𝑛 (𝑛 = 1,2, … ) is a sequence of pairwise exclusive events and 𝐸𝑛 , we
have

𝑃 𝐸 = 𝑃 𝐸𝑛
𝑛=1
This property is called the additivity of probability. If 𝑃 𝐸 = 1, the event 𝐸 is said to
occur almost certainly (almost surely (abbrev. a.s.), for almost all 𝜔 or with probability
1).

PRODUCT MEASURE: Let 𝑋 and 𝑌 be two spaces, and let 𝑆 and 𝑇 be semirings on 𝑋 and 𝑌
respectively. Then 𝑆 × 𝑇 is the semiring consisting of { 𝐴 × 𝐵 ∶ 𝐴 ∇ 𝑆, 𝐵 ∇ 𝑇 }. Let µ
and 𝜈 be measures on 𝑆 and 𝑇 respectively. Define the product measure µ × 𝜈 on 𝑆 × 𝑇
by the rule (µ × 𝜈)(𝐴 × 𝐵) = µ(𝐴) 𝜈(𝐵).

PROJECTION: Our shadow is our projection on the ground. x is the projection on the X
axis and y is the projection on the Y axis, of the point (𝑥, 𝑦).

PROJECTION THEOREM: Let 𝑋 be a Hilbert space, 𝐾 a closed convex subset, and 𝑥 ∇ 𝑋.

Then there exists a unique 𝑥 ∇ 𝐾 such that 𝑘𝑥 − 𝑥𝑘 = 𝑖𝑛𝑓 𝑦 ∇ 𝐾 𝑘𝑥 − 𝑦𝑘.

PROPER IDEAL: An ideal 𝐼 of 𝑅 is said to be a proper ideal of 𝑅 if 𝐼 ≠ 𝑅.

Proportionality: In mathematics, proportionality indicates that two quantities or

variables are related in a linear manner. If one quantity doubles in size, so does the
other; if one of the variables diminishes to 1/10 of its former value, so does the other.
The symbol for proportionality resembles a stretched-out, lowercase Greek letter alpha
( ). When this symbol appears between two quantities or variables, it is read "is
proportional to" or "varies in direct proportion with." Thus, the expression x y is
read " x is proportional to y " or " x varies in direct proportion with y ." In this situation,
as long as x and y do not attain values of zero, the quotient x / y is always equal to the
same value k , which is called the proportionality constant.

PSEUDOCONCAVE FUNCTION: Let θ be a numberical function defined on some open set

in Rn containing the set 𝑇.

θ is said to be pseudoconcave at x ∇ T (with respect to T), if it is differentiable at x and

X ∇ T, ∆θ x x − x ≥ 0 ⇒ θ x ≤ 𝜃(𝑥).

θ is said to be pseudiconcave on 𝑇 if it is pseudoconcave at each x ∇ T.

Clearly θ is pseudoconcave at x(on T) if and only if – θ is pseudiconvex at x(on T).

PSEUDOCONVEX FUNCTION: Let θ be a numberical function defined on some open set in
Rn containing the set 𝑇.

θ is said to be pseudoconvex at x ∇ T (with respect to T), if it is differentiable at x and

x ∇ T, ∆θ x x − x ≥ 0 ⇒ θ x ≥ 𝜃(𝑥).

θ is said to be pseudiconvex on 𝑇 if it is pseudoconvex at each x ∇ T.

PURE SHEAR: If a volume initially cubical is deformed into a parallelepiped and the
volumes of the cube and the parallelepiped are equal, when the product of changes is
neglected in the linear elements, the deformation is called Pure-shear.

PYTHAGORAS THEOREM FOR ORTHOGONAL VECTORS: If 𝑥 ⊥ 𝑦 then ||𝑥 + 𝑦||2 =

||𝑥||2 + ||𝑦||2 .

QED: The term “QED” is actually an abbreviation and stands for the Latin quod erat
demonstrandum, meaning “which was to be demonstrated.” QED typically is used to
signify the end of a mathematical proof.
QUADRATIC DUAL (MAXIMIZATION) PROBLEM (QDP): The quadric dual
(maximization) problem means to find an x ∇ Rn and a u ∇ Rm , if they exist, such that

1 1
− 𝑥 𝐶 𝑥 − 𝑐 𝑢 = max − 𝑥 𝐶 𝑥 − 𝑐𝑢
2 (𝑥,𝑢)∇𝑌 2

where x, u ∇ 𝑌 = x, u : x ∇ Rn , u ∇ Rm , Cx + A′ u = b and u ≥ 0

QUADRATIC FORM: An expression of the form

n n

aij xi xj ,
i−1 j−1

where aij ’s elements of a field F, is called a quadratic form in the n variables

x1 , x2 , … … , xn over a field F.

QUADRATIC FORM CORRESPONDING TO A SYMMETRIC MATRIX: Let A = aij be a

n×n

symmetrix matrix over the field 𝐹 and let 𝑋 = [x1 , x2 , … … , xn ]𝑇 be a column vector.
n n
Then 𝑋 𝑇 𝐴𝑋 determines a unique quadratic form i−1 j−1 a ij xi xj in n variables
x1 , x2 , … … , xn over the field F.
QUADRIC HYPERSURFACES: A subset 𝐹 of a n-dimensional Euclidean space 𝐸 𝑛 is called a
quadric hypersurface (or simply hyperquadric) if it is the set of points 𝑥1 , … , 𝑥𝑛
satisfying the following equation of the second degree:

𝑛 𝑛
𝑖,𝑘=1 𝑎𝑖𝑘 𝑥𝑖 𝑥𝑘 +2 𝑖=1 𝑏𝑖 𝑥𝑖 + 𝑐 = 0,

where 𝑎𝑖𝑘 , 𝑏𝑖 , 𝑐 are all real numbers.

QUADRATIC FIELD: A quadratic field is an algebraic number field of degree 2.

QUADRATIC PROGRAMMING: A quadratic programming problem is a special type of

mathematical programming where the objective function is quadratic while the
constraints are linear. A typical formulation of the problem is as follows.

(𝑄) Maximize 𝑧 = 𝑐 ′ 𝑥 − 12 𝑥′𝐷𝑥 under the condition 𝐴𝑥 ≤ 𝑏 𝑎𝑛𝑑 𝑥 ≥ 0, 𝑥 ∇ 𝑅 𝑛 .

QUADRIC SURFACES: A subset 𝐹 of a 3-dimensional Euclidean space 𝐸 3 is called a

quadric surface (surface of the second order or simply quadric) if 𝐹 is the set of zeroes
of a quadratic equation 𝐺 𝑥, 𝑦, 𝑧 = 0, where the coefficients of 𝐺 are real numbers. The
equations 𝐺 𝑥, 𝑦, 𝑧 = 0 is written as 𝑎𝑥 2 + 𝑏𝑦 2 + 𝑐𝑥 2 + 𝑑 + 2𝑓𝑦𝑧 + 2𝑔𝑧𝑥 + 2𝑕𝑥𝑦 +
2𝑓 ′ 𝑧 + 2𝑔′ 𝑦 + 2𝑕′ 𝑧 = 0.

QUADRATURE: Quadrature is the computation of a univariate definite integral. It can

refer to either numerical or analytic techniques; one must gather from context which is
meant. Cubature refers to higher-dimensional definite integral computation. Some
numerical quadrature methods are Simpson’s rule, the trapezoidal rule, and Riemann
sums.
QUASICONCAVE FUNCTION: A numerical function θ defined on a set T ⊂ Rn is said to be
quasiconcave at x ∇ T (with respect to T) if for each x ∇ T such that θ x ≥ θ(x), the
function θ assumes a values no larger smaller then θ(x) on each point in the intersection
of the closed line segment [x, x] and 𝑇 , or equivalently

X ∇ T, θ x ≤ x , 0 ≤ λ ≤ 1 & 1 − 𝜆 x + λx ∇ T ⇒ θ(x) ≤ 𝜃 1 − 𝜆 x + λx

𝜃 is said to be quasiconcave on 𝑇 if it is quasiconcave at each x ∇ T.

QUASICONFORMAL MAPPINGS: H.Grotzsh (1928) introduced quasiconformal mappings

as a generalization of conformal mappings. Let 𝑓(𝑧) be a continuously differentiable
homeomorphism with positive Jacobian between place domains. This image of an
infinitesimal circle 𝑑𝑧 = constant is an infinitesimal ellipse with major axis of length
𝑓𝑧 + 𝑓𝑧 𝑑𝑧 and minor axis of length 𝑓𝑧 + 𝑓𝑧 𝑑𝑧 . When the ratio 𝐾 𝑧 =
𝑓𝑧 + 𝑓𝑧 𝑓𝑧 − 𝑓𝑧 is bounded, 𝑓is called quasiconformal. If 𝐾 = 1, then 𝑓 is
conformal.

QUASICONVEX FUNCTIONS: A numerical function θ defined on a set T ⊂ Rn is said to be

quasiconvex at x ∇ T (with respect to T) if for each x ∇ T such that θ x ≤ θ(x), the
function θ assumes a values no larger then θ(x) on each point in the intersection of the
closed line segment [x, x] and 𝑇 , or equivalently

x ∇ T, θ x ≤ x , 0 ≤ λ ≤ 1 & 1 − 𝜆 x + λx ∇ T ⇒ 𝜃 1 − 𝜆 x + λx ≤ θ(x)

𝜃 is said to be quasi convex on 𝑇 if it is quasiconvex at each x ∇ T.

QUASI FUZZY NUMBER: A quasi fuzzy number 𝐴 is a fuzzy set of the real line with a
normal, fuzzy convex and continuous membership function satisfying the limit
conditions
lim 𝐴 𝑡 = 0, lim 𝐴 𝑡 = 0
𝑡→∞ 𝑡→−∞

QUATERNIONS: A quaternion may be defined to be an expression of the form

𝑤 + 𝑥𝑖 + 𝑦𝑗 + 𝑧𝑘, where 𝑤, 𝑥, 𝑦 and 𝑧 are real numbers. There are operations of
addition, subtraction and multiplication defined on the set 𝐻 of quaternions. These are
binary operations on that set. Quaternions were introduced into mathematics in 1843
by William Rowan Hamilton. The definitions of addition and subtraction are
straightforward. The sum and difference of two quaternions 𝑤 + 𝑥𝑖 + 𝑦𝑗 + 𝑧𝑘 and
𝑤 ′ + 𝑥 ′ 𝑖 + 𝑦 ′ 𝑗 + 𝑧 ′ 𝑘 are given by the formulae

𝑤 + 𝑥𝑖 + 𝑦𝑗 + 𝑧𝑘 + 𝑤 ′ + 𝑥 ′ 𝑖 + 𝑦 ′ 𝑗 + 𝑧 ′ 𝑘 = 𝑤 + 𝑤 ′ + 𝑥 + 𝑥 ′ 𝑖 +
𝑦 + 𝑦 ′ 𝑗 + 𝑧 + 𝑧 ′ 𝑘;

(𝑤 + 𝑥𝑖 + 𝑦𝑗 + 𝑧𝑘) − (𝑤 ′ + 𝑥 ′ 𝑖 + 𝑦 ′ 𝑗 + 𝑧 ′ 𝑘) = (𝑤 − 𝑤 ′ ) + (𝑥 − 𝑥 ′ )𝑖 +
(𝑦 − 𝑦 ′ )𝑗 + (𝑧 − 𝑧 ′ )𝑘.

If the quaternions 𝑤 + 𝑥𝑖 + 𝑦𝑗 + 𝑧𝑘 and 𝑤 ′ + 𝑥 ′ 𝑖 + 𝑦 ′ 𝑗 + 𝑧 ′ 𝑘 are denoted by 𝑞

and 𝑞 ′ respectively, then we may denote the sum and the difference of these
quaternions by 𝑞 + 𝑞 ′ and 𝑞 − 𝑞 ′ . These operations of addition and subtraction of
quaternions are binary operations on the set 𝐻 of quaternions. It is easy to see that the
operation of addition is commutative and associative, and that the zero quaternion
0 + 0𝑖 + 0𝑗 + 0𝑘 is an identity element for the operation of addition. The operation
of subtraction of quaternions is neither commutative nor associative. This results
directly from the fact that the operation of subtraction on the set of real numbers is
neither commutative nor associative.

QUEUE LENGTH: Number of persons in the system at any time. We may be interested in
studying the distribution of queue length.

QUEUING MODEL (MARKOV CHAINS H): In a telephone system, calls made when all the
lines of the system are busy are lost. The problem of computing the probability of loss
involved was fïrst solved in the pioneering article on queuing theory by A. K. Erlang in
1917. For systems in which calls can be put on hold when all lines are busy, one deals
with the waiting time distribution instead of the probability. In 1930, F. Pollaczeck and
A. Ya. Khinchin gave explicit formulas for the characteristic function of the waiting time
distribution. In many situations, such as in telephone systems, air and surface traffic,
production lines, and computer systems, various congestion phenomena are often
observed, and many kinds of queuing models are utilized to analyze the congestion.
Mathematically, almost all such models are formulated by using Markov processes. For
practical uses, approximation and computational methods are important as well as
theoretical results.
QUEUING PROCESS: A queuing process is centered on a service system (facility). The
queueing system has the following characteristics:

(i) The input( arrival pattern)

(ii) Queue (waiting line)
(iii) The service discipline (queue discipline)
(iv) The service mechanism( service pattern)

All queuing situation involve the arrival of customers (input) at a service facility, where
some time may be spent in waiting and then receiving the desired service. The
customers arriving for service may or may not enter the system. Thus, the input pattern
in the system depends on the nature of the system as well as the behavior of the
customer. The combination of these two determines the arrival pattern. After the
service is completed the customer leaves the service system (output). The departure
pattern mainly depends on the service discipline. There can be many types of queuing
systems depending on the nature of inputs, service mechanism and customer
characteristics.

QUEUING THEORY: A group of customers / items, waiting at some place to receive

attention/ service including those receiving the service, is known as queue.

Similarly some service facility waits for arrival of customers when the total capacity of
system is more than the number of customers, Thus , in the absence of a perfect balance
between the service facility and the customers, waiting is required either by the service
facility or by the customer. The imbalance between the customers and service facility ,
known as congestion , cannot be eliminated completely but efforts/ techniques can be
evolved to reduce the magnitude of congestion or waiting time of a new arrival in the
system.

The customers arriving for service may from one queue and be serviced through only
one station , that may form one queue and be serviced through several stations or they
may form several queues and be served through as many stations.

QUOTIENT LAW: A set of quantities, whose inner product with an arbitrary vector is a
tensor, it itself a tensor.

QUOTIENT MAP: Let 𝑋 and 𝑌 be two topological spaces and 𝑝 ∶ 𝑋 → 𝑌 be a surjective

map. The map 𝑝 is said to be a quotient map provided a subset 𝑈 of 𝑌 is open in 𝑌 if and
only if 𝑝−1 (𝑈) is open in 𝑋.
QUOTIENT MODULE: If 𝑁 is a submodule of a module 𝑀 over a ring 𝑅 we write 𝑀/𝑁 for
the set of cosets of 𝑁 and define addition and multiplication on 𝑀/𝑁 by (𝑢 + 𝑁) +
(𝑣 + 𝑁) = (𝑢 + 𝑣) + 𝑁, 𝑟(𝑢 + 𝑁) = 𝑟𝑢 + 𝑁. If 𝑁 is a submodule of a module 𝑀
over a ring 𝑅 then 𝑀/𝑁 with module addition and multiplication is a module over 𝑅. We
call 𝑀/𝑁 a quotient module.
QUOTIENT RING: If 𝐼 is an ideal of a ring 𝑅 we write 𝑅/𝐼 for the set of cosets of 𝐼 and
define addition and multiplication on 𝑅/𝐼 by (𝑟 + 𝐼) + (𝑠 + 𝐼) = (𝑟 + 𝑠) + 𝐼, (𝑟 +
𝐼)(𝑠 + 𝐼) = 𝑟𝑠 + 𝐼. If 𝐼 is an ideal of a ring 𝑅 then 𝑅/𝐼 with addition and
multiplication is a ring. We call 𝑅/𝐼 a quotient ring.
RAABE’S TEST: The series 𝑢𝑛 of positive terms is convergent or divergent according as
𝑢𝑛
lim𝑛→∞ 𝑛 −1 > 1 𝑜𝑟 < 1.
𝑢 𝑛 +1

The radian is the Standard International (SI) unit of plane angular measure. There are 2 pi, or
approximately 6.28318, radians in a complete circle. Thus, one radian is about 57.296
angular degrees.

RADIAN: The term radian arises from the fact that the length of a circular arc,
corresponding to an angle of one radian, is equal to the radius of the arc. This is shown
in the illustration. Point 𝑃 represents the center of the circle. The angle 𝑞, representing
one radian, is such that the length of the subtended circular arc is equal to the radius, 𝑟,
of the circle.

The radian is used by mathematicians, physicists, and engineers. It arises in natural

phenomena and in equations which, unlike the angular degree, were invented for
human convenience.

RADICAL PLANE: The radical plane of any two given spheres is the locus of a point
which moves so that the squares of the lengths of the tangents drawn from it to the
given spheres are equal.

RADICALS: let 𝐴 be a ring. Then among ideals consisting only of quasi-invertible

elements of 𝐴, there exists a largest one, which is called the radical of 𝐴 and denoted by
ℜ(𝐴). The radical of the residue ring 𝐴/ℜ(𝐴) is {0}. A ring 𝐴 is called a semiprimitive
ring if ℜ(𝐴) is {0}. On the other hand, 𝐴 is called a left (right) primitive ring if it has a
faithful simple left(right) 𝐴-module. The radical ℜ(𝐴) is equal to the intersection of all
ideals 𝐽 such that 𝐴/𝐽 is a left (right) primitive ring. In a unitary ring 𝐴, ℜ(𝐴) coincides
with the intersection of all maximal left (right) ideal of 𝐴. Furthermore, in a left (right)
Artinian ring 𝐴, ℜ(𝐴) is the largest nilpotent ideal of 𝐴, and the condition 𝑅 𝐴 = {0} is
equivalent to the condition that 𝐴 is a semi-simple ring.

∞ 𝑛
RADIUS OF CONVERGENCE: Let 𝑎𝑛 ∇ 𝐶. Then, either 𝑛=0 𝑎𝑛 𝑧 converges for all 𝑧 ∇ 𝐶
(and we say that the power series has radius of covergence infinity) or there exists an 𝑅
with 𝑅 ≥ 0 such that
∞ 𝑛
(i) 𝑛=0 𝑎𝑛 𝑧 converges for all |𝑧| < 𝑅,
∞ 𝑛
(ii) 𝑛=0 𝑎𝑛 𝑧 diverges for all |𝑧| > 𝑅.

We call R the radius of convergence.

RADON'S THEOREM: Any set of 𝑑 + 2 points in 𝑹𝒅 can be partitioned into two disjoint
sets whose convex hulls intersect. A point in the intersection of these hulls is called
a Radon point of the set.

RADON–NIKODYM THEOREM: Any charge 𝜈 which absolutely continuous with respect

to a measure µ have the form ν(A) = ∫𝐴 f dµ, where f is a function from 𝐿1 . The function
𝑓 ∇ 𝐿1 is uniquely defined by the charge 𝜈.

Let µ be a measure on 𝑋, ν be a finite charge, which is absolutely continuous with

respect to µ. For any 𝜀 > 0 there exists 𝛿 > 0 such that µ(𝐴) < 𝛿 implies | 𝜈 |(𝐴) < 𝜀 .

RANDOM NUMBERS: A sequence of numbers that can be regarded as realizations of

independent and identically distributed random variables is called a sequences or table
of random numbers. It is a basic tool for the Monte Carlo method, simulation of
stochastic phenomena in nature or in society, and sampling or randomization
techniques in statistics. Random numbers are numbers that occur in a sequence such
that two conditions are met:

(1) the values are uniformly distributed over a defined interval or set, and

(2) it is impossible to predict future values based on past or present ones. Random
numbers are important in statistical analysis and probability theory.
The most common set from which random numbers are derived is the set of single-digit
decimal numbers {0, 1, 2, 3, 4, 5, 6, 7, 8, 9}. The task of generating random digits from
this set is not trivial. A common scheme is the selection (by means of a mechanical
escape hatch that lets one ball out at a time) of numbered ping-pong balls from a set of
10, one bearing each digit, as the balls are blown about in a container by forced-air jets.
This method is popular in lotteries. After each number is selected, the ball with that
number is returned to the set, the balls are allowed to blow around for a minute or two,
and then another ball is allowed to escape.

Sometimes the digits in the decimal expansions of irrational numbers are used in an
attempt to obtain random numbers. Most whole numbers have irrational square roots,
so entering a string of six or eight digits into a calculator and then hitting the square
root button can provide a sequence of digits that seems random. Other algorithms have
been devised that supposedly generate random numbers. The problem with these
methods is that they violate condition (2) in the definition of randomness. The existence
of any number-generation algorithm produces future values based on past and/or
current ones. Digits or numbers generated in this manner are called pseudorandom.

Statisticians, mathematicians, and scientists have long searched for the ideal source of
random numbers. One of the best methods is the sampling of electromagnetic noise.
This noise, generated by the chaotic movements of electrons, holes, or other charge
carriers in materials and in space, is thought to be as close to "totally random" as any
observable phenomenon.

RANGE OF A FUNCTION: Let 𝐴 and 𝐵 be sets, and let 𝑓: 𝐴 → 𝐵 be a function from 𝐴 to 𝐵.

The range of the function f is the subset 𝑓(𝐴) of 𝐵 defined by 𝑓(𝐴) = {𝑏 ∇ 𝐵 ∶ 𝑏 =
𝑓(𝑎) for some 𝑎 ∇ 𝐴}. In other words, the range of a function is the set consisting of all
elements of the codomain of the function that are images under the function of elements
of its domain.

RANKINE’S COMBINED VORTEX: Rankine’s combined vortex consists of a circular

cylindrical vortex with its axis vertical in a liquid moving irrotationally under the action
of gravity, only the upper surface being exposed to atmospheric pressure Π. The
external forces are derivable from the derivable form the potential gz.
RANK-MULTIPLICITY THEOREM: The geometric multiplicity of a characteristic root
cannot exceed its algebraic multiplicity.

RANK-NULLITY THEOREM: Let 𝑈, 𝑉 be vector spaces over 𝐾 with 𝑈 finite-dimensional

and let 𝑇: 𝑈 → 𝑉 be a linear map.Then

𝑟𝑎𝑛𝑘(𝑇) + 𝑛𝑢𝑙𝑙𝑖𝑡𝑦 (𝑇) = 𝑑𝑖𝑚 𝑈

RATE OF A CODE: Let 𝐶 be an (𝑛, 𝑀)-code over an alphabet of size 𝑞. Then, the rate of
the code 𝐶 is defined by
log 𝑞 𝑀
𝑟𝑎𝑡𝑒(𝐶) =
𝑛
RATIONAL NUMBER: A rational number is a number determined by the ratio of
some integer p to some nonzero natural number q. The set of rational numbers is
denoted Q, and represents the set of all possible integer-to-natural-number ratios p/q.
In mathematical expressions, unknown or unspecified rational numbers are
represented by lowercase, italicized letters from the late middle or end of the alphabet,
especially r, s, and t, and occasionally u through z. Rational numbers are primarily of
interest to theoreticians. Theoretical mathematics has potentially far-reaching
applications in communications and computer science, especially in data encryption and
security.

If r and t are rational numbers such that r < t, then there exists a rational number s such
that r < s < t. This is true no matter how small the difference between r and t, as long as
the two are not equal. In this sense, the set Q is "dense."Nevertheless, Q is
a denumerable set. Denumerability refers to the fact that, even though a set might
contain an infinite number of elements, and even though those elements might be
"densely packed," the elements can be defined by a list that assigns them each a unique
number in a sequence corresponding to the set of natural numbers N = {1, 2, 3, ...}..

For the set of natural numbers N and the set of integers Z, neither of which are "dense"
denumeration lists are straightforward. For Q, it is less obvious how such a list might be
constructed. An example appears below. The matrix includes all possible numbers of the
form p/q, where p is an integer and q is a nonzero natural number. Every possible
rational number is represented in the array. Following the pink line, think of 0 as the
"first stop," 1/1 as the "second stop," -1/1 as the "third stop," 1/2 as the "fourth stop,"
and so on. This defines a sequential (although redundant) list of the rational numbers.
There is a one-to-one correspondence between the elements of the array and the set of
natural numbers N.

To demonstrate a true one-to-one correspondence between Q and N, a modification

must be added to the algorithm shown in the illustration. Some of the elements in the
matrix are repetitions of previous numerical values. For example, 2/4 = 3/6 = 4/8 =
5/10, and so on. These redundancies can be eliminated by imposing the constraint, "If a
number represents a value previously encountered, skip over it."In this manner, it can
be rigorously proven that the set Q has exactly the same number of elements as the
set N. Some people find this hard to believe, but the logic is sound.

In contrast to the natural numbers, integers, and rational numbers, the sets of irrational
numbers, real numbers, imaginary numbers, and complex numbers are non-
denumerable. They have cardinality greater than that of the set N.

RANK OF A FREE MODULE: The rank of a free module is the number of elements in any
free basis for the free module.

RANK OF A MATRIX: A number is said to be the rank of a matrix 𝐴 if it possesses the

following two properties:

(i) There is at least one square submatrix of 𝐴 of order 𝑟 whoes determinant is not
equal to zero.
(ii) If the matrix 𝐴 contains any square submatrix of order 𝑟 + 1, then the
determinant of every square submatrix of 𝐴 of order 𝑟 + 1 should be zero.

In short the rank of a matrix is the order of any highest order non-vanishing minor if 𝑡𝑕
matrix.

Since the rank of every non-zero matrix is ≥ 1, we agree to assign the rank, zero, to
every null matrix.

(i) The rank of a matrix is ≤ r, if all 𝑟 + 1-rowed minors of the matrix vanish.
(ii) The rank of a matrix is ≥ r, if there is at least one 𝑟-rowed minors of the matrix
which is not equal to zero.

RANK OF A QUADRATIC FORM: Let X ′ AX be a quadratic form over a field F . The rank of
the matrix A is called the rank of the quadratic form X ′ AX.

RANK OF A SUM: Rank of the sum of two matrices cannot exceed sum of their ranks.

RATE OF TRANSMISSION OF ENERGY: The rate of transmission of energy is measured by

taking a vertical section of the liquid at right angles to the direction of propagation. We
shall determine the rate at which the liquid on one side of this section is doing work on
the liquid on the other side.

RATIO TEST: If 𝑎𝑛 ∇ 𝑅, 0 < 𝑎𝑛 and 𝑎𝑛+1 /𝑎𝑛 → 𝑙 then

∞
(iii) If 𝑙 < 1 then 𝑛=1 𝑎𝑛 converges.
∞
(iv) If 𝑙 > 1 then 𝑛=1 𝑎𝑛 diverges.

REAL NUMBER: A real number is any element of the set R, which is the union of the set
of rational numbers and the set of irrational numbers. In mathematical expressions,
unknown or unspecified real numbers are usually represented by lowercase italic
letters u through z. The set R gives rise to other sets such as the set of imaginary
numbers and the set of complex numbers. The idea of a real number (and what makes it
"real") is primarily of interest to theoreticians. Abstract mathematics has potentially
far-reaching applications in communications and computer science, especially in data
encryption and security.
If x and z are real numbers such that x < z, then there always exists a real
number y such that x < y < z. The set of reals is "dense" in the same sense as the set of
irrationals. Both sets are nonde numerable. There are more real numbers than is
possible to list, even by implication.

The set R is sometimes called the continuum because it is intuitive to think of the
elements of R as corresponding one-to-one with the points on a geometric line. This
notion, first proposed by Georg Cantor who also noted the difference between the
cardinalities (sizes) of the sets of rational and irrational numbers, is called
the Continuum Hypothesis. This hypothesis can be either affirmed or denied without
causing contradictions in theoretical mathematics.

REAL OR ACTUAL FLUIDS: The real fluid (Viscous and compressible) is one in which
both the tangential and normal forces exist. Viscosity, which is also known as an
internal friction, of a fluid is that characteristic of real fluid which is capable to offer
resistance to shearing stress. The resistance is comparatively small (not negligible) for
fluids such as water and gases but it quite large for other fluids such as oil, glycerin,
molasses, coal tar etc.

REAL QUADRATIC FORM: An expression of the form

n n

aij xi xj ,
i−1 j−1

Where aij ’s are all real numbers, is called a real quadratic form in the n variables
x1 , x2 , … … , xn For example,

(i) 2x 2 + 7xy + 5y 2 is a real quadratic form in the two variables x and y.

(ii) 2x 2 − y 2 + 2z 2 − 2yz − 4zx + 6xy is a real quadratic form in the three variables
x, y and z.
(iii) x12 − 2x22 + 4x32 − 4x42 − 2x1 x2 + 3x1 x4 + 4x2 x3 − 5x3 x4 is a real quadratic form
in the four variables x1 , x2 , x3 and x4 .

REAL VECTOR SPACES: A real vector space consists of a set 𝑉 on which there is defined
an operation of vector addition, yielding an element 𝑣 + 𝑤 of 𝑉 for each pair 𝑣, 𝑤 of
elements of 𝑉 , and an operation of multiplication-by-scalars that yields an element 𝜆𝑣
of 𝑉 for each 𝑣 ∇ 𝑉 and for each real number 𝜆. The operation of vector addition is
required to be commutative and associative. There must exist a zero element 0𝑣 of V
that satisfies 𝑣 + 0𝑣 = 𝑣 for all 𝑣 ∇ 𝑉 , and, for each 𝑣 ∇ 𝑉 there must exist an
element −𝑣 of 𝑉 for which 𝑣 + (−𝑣) = 0𝑣 . The following identities must also be
satisfied for all 𝑣, 𝑤 ∇ 𝑉 and for all real numbers 𝜆 and µ: (𝜆 + µ)𝑣 = 𝜆𝑣 + µ𝑣, 𝜆(𝑣 +
𝑤) = 𝜆𝑣 + 𝜆𝑤, 𝜆(µ𝑣) = (𝜆µ)𝑣, 1𝑣 = 𝑣.

RECURRENCE FORMULAE FOR ASSOCIATED LEGENDRE FUNCTION:

2mx
(i) Pn m+1 x − Pn m x + n n + 1 − (m − 1) Pn m−1 x = 0.
1−x 2

(ii) 2n + 1 = xP m n x = n + m − P m n−1 x − (n − m + 1)P m n+1 x .

1 m+1 m+1
(iii) 1 − x 2 P m n (x) = (2n+1) Pn+1 x − Pn−1 x .
1 m−1
(iv) 1 − x 2 P m n x + (2n+1) n + m n + m − 1 Pn−1 x − n − m + 1 (n −
m−1
m − 2)Pn+1 x .

RECURRENCE FORMULAE FOR HERMITE POLYNOMIALS:

(i) Hn ′ x = 2nHn−1 x n ≥ 1.
(ii) 2x. Hn x = 2nHn−1 x + Hn+1 (x)
(iii) H′n x = 2xHn x − Hn+1 (x).
(iv) Hn " x − 2x H′n x + 2nH x = 0.

RECURRENCE FORMULAE FOR 𝐉𝐧 (𝐱):

(i) x. J ′ n x = n. Jn x − x. Jn+1 (x)

(ii) x. J ′ n x = −nJn x + xJn−1 (x)
(iii) 2 J ′ n x = Jn−1 x − Jn+1 (x).
(iv) 2nJn x = x[Jn−1 x + Jn+1 x ]
d
(v) [x −n Jn x ] = x −n Jn+1 x .
dx

(vi) d/dx[x n Jn x ] = x n Jn+1 x .

RECURRENCE FORMULAE FOR LAGUERRE POLYNOMIALS:

(i) n + 1 Ln+1 x = 2n + 1 − x nLn x − Ln−1 (x).

(ii) xL′n x = nLn x + nLn−1 (x)
 n−1
(iii) L′n x = − r=0 Lr (x).

RECURRENCE FORMULAE FOR LEGENDRE’S EQUATION:

(i) 2n + 1 x𝑃𝑛 = n + 1 Pn+1 + nPn−1 .

(ii) nPn = xP′n − P′n−1, where dashes denotes differential w. r. t.′ x ′ .
(iii) 2n + 1 Pn = P′n+1 − P′n−1 .
(iv) n + 1 Pn = P′n+1 − xP′n .
(v) 1 − x 2 P′n = n Pn+1 − xPn .
(vi) 1 − x 2 P′n = n + 1 xPn − Pn+1 .

RECURRENCE FORMULAE FOR 𝐐𝐧 (𝐱):

(i) Q′n+1 x − Q′n−1 x = (2n + 1)Qn x .

(ii) Q′n+1 x − xQ′n x = (n + 1)Qn x .
(iii) nQ′n+1 x + (n + 1)Q′ n−1 x = (2n + 1)xQn ′ x .
(iv) xQn ′ x − Q′n−1 x = nQn x .
(v) x 2 − 1 Q′n x = nxQn x − nQn−1 x .
(vi) n + 1 Qn+1 x − 2n + 1 xQn x + nQn−1 x = 0.

RECURRENCE FORMULAE FOR 𝐓𝐧 𝐱 AND 𝐔𝐧 𝐱 :

(i) Tn+1 x − 2x Tn x + Tn−1 x = 0.

(ii) (1 − x 2 )T′n x = −nx Tn x + nTn−1 x
(iii) Un+1 x − 2x Un x + Un−1 x = 0
(iv) 1 − x 2 U ′ n x = −nx Un x + nUn−1 x

RECIPROCAL SYSTEM OF VECTORS: If 𝒂, 𝒃, 𝒄 be any three non-coplanar vectors so that

𝒂𝒃𝒄 ≠ 0 then the three vectors 𝒂’, 𝒃’, 𝒄’ defined by the equations

𝒃×𝒄 𝒄×𝒂 𝒂×𝒃

𝒂′ = ,𝒃 ′ = , 𝒄′ =
𝒂𝒃𝒄 𝒂𝒃𝒄 𝒂𝒃𝒄

are called the reciprocal system of vectors to the vectors 𝒂, 𝒃, 𝒄.

RECTIFYING PLANE: The plane containing the tangent and binormal at 𝑃 is called
rectifying plane at 𝑃. i.e., it is plane passing through 𝑃 and perpendicular to principal
normal at 𝑃. Evidently equation of this plane is
 𝑅−𝑟 ∙𝑛 =0
RECTILINEAR VORTICES: Consider a two dimensional vortex motion. The vorticity
vector is perpendicular to the plane of motion and the vortex lines are straight lines
parallel to 𝑍-axis. All vortex tubes are cylindrical, with generators perpendicular to the
plane of the motion. Such vortices are called rectilinear vortices.

REDUCIBILITY TOPOLOGICAL SPACE: A topological space 𝑍 is said to be reducible if it

can be decomposed as a union 𝐹1 ∪ 𝐹2 of two proper closed subsets 𝐹1 and 𝐹2 . A
topological space 𝑍 is said to be irreducible if it cannot be decomposed as a union of two
proper closed subsets.

REDUCTION TO NORMAL FORM: Every 𝑚 × 𝑛 matrix of the rank 𝑟can be reduced to the
Ir 0
form by a finite chain of 𝐸- operations, where Ir is the 𝑟- rowed unit matrix.
0 0

Equivalence of Matrices: If 𝐵 be any 𝑚 × 𝑛 matrix obtained from an 𝑚 × 𝑛 matrix 𝐴 by

finite number of elementary transformations of 𝐴, then 𝐴 is called equivalent to 𝐵.
Symbolically, we write 𝐴~𝐵, which is read as ‘𝐴is equivalent to 𝐵’.

The following three properties of the relation ‘~’ in the set of all 𝑚 × 𝑛 matrices are
quite obvious:

(i) Reflexivity: If 𝐴 is any 𝑚 × 𝑛 matrix, then 𝐴~𝐴. Obviously 𝐴 can be obtained form
𝐴 by the elementary transformation R i → kR i , where k = 1.
(ii) Symmetry: If 𝐴~𝐵, then 𝐵~𝐴. If 𝐵 can be obtained from 𝐴 by a finite number of
elementary transformation of 𝐴, then 𝐴 can also be obtained from 𝐵 by a finite number
of elementary transformations of 𝐵.
(iii) Transitivity: If 𝐴~𝐵, then 𝐵~𝐶, and 𝐴~𝐶.
If 𝐵 can be obtained from 𝐴 by a series of elementary transformation of 𝐴 and 𝐶 can be
obtained from 𝐵 by a series of elementary transformations of 𝐵, then 𝐶 can also be
obtained from 𝐴 by a series of elementary transformations of 𝐴.

Therefore the relation ‘~’ in the set of all 𝑚 × 𝑛 matrices is equivalence relations.

REFINEMENT OF A PARTITION: Let 𝑃 and 𝑃 ∗ be two partitions of a closed interval [𝑎, 𝑏].
Then 𝑃∗ is said to be a refinement of 𝑃 if 𝑃 ⊆ 𝑃 ∗ . We also say that 𝑃 ∗ is finer than 𝑃.
REGULAR ELEMENTS: Let A be Banach algebra with identity. If 𝑓 is an element in A,
then it may happen that there is an element 𝑔 in A such that 𝑓𝑔 = 𝑔𝑓 = 𝑒 (identity). In
this case 𝑔 is denoted by 𝑓 −1 and called to inverse of 𝑓.
If an element 𝑓 in has an inverse, then 𝑓 is said to be regular element. Regular elements
are often called invertible elements or non-singular elements or unit element. The
element 𝑒 is always regular in a Banach algebra A. the set of regular elements in A is
denoted by G. Clearly G contains 𝑒 and is a group.
REGULAR POLYGONS: A ‘polygon in a Euclidean plane bounding a convex cell whose
sides and interior angles are all respectively congruent is called a regular polygon.
When the number of vertices (which equals the numbers of sides) is 𝑛, it is called a
regular 𝑛 −gon. There exist a circle (circumscribed circle) passing through all the
vertices of a regular 𝑛- gon and a concentric circle (inscribed circle) tangent to all the
sides. We call the center of those circles the center of the regular 𝑛-gon. The 𝑛 vertices
of a regular 𝑛- gon are obtained by divided a circle into 𝑛 equal parts.

REGULAR SPACES: The space 𝑋 is regular if any 𝑥 ∇ 𝑋 and 𝐶 ⊂ 𝑋, with 𝐶 closed

and 𝑥 not in 𝐶, have disjoint neighborhoods. 𝑋 is regular if the neighborhood filter of
each point has a base consisting of closed sets.

REGULAR VALUE (MANIFOLD): An element 𝑦 ∇ 𝑁 is called a regular value if for each

𝑥 ∇ 𝑓 −1 (𝑦) = {𝑥 ∇ 𝑀 | 𝑓(𝑥) = 𝑦}, the differential 𝑑𝑓𝑥 : 𝑇𝑥 𝑀 → 𝑇𝑦 𝑁 is surjective. A
critical value is an element 𝑦 ∇ 𝑁 which is not regular. Notice that by this definition all
elements 𝑦 ∈ 𝑓(𝑀) are regular values. Notice also that if there exists a regular value
𝑦 ∇ 𝑓(𝑀), then necessarily 𝑛 ≤ 𝑚, since a linear map into a higher dimensional space
cannot be surjective.
RELATION BETWEEN BETA AND GAMMA FUNCTIONS: Relation between Beta and
Gamma Functions is given as

(Γ(𝑙)Γ(𝑚 ))
𝐵 𝑙, 𝑚 = (Γ(𝑙+𝑚 ))

RELATIONS: In its wider sense the term relation means 𝑛- ary relation (𝑛 = 1,2,3, … )
but we restrict ourselves to its most ordinary meaning, i.e. ,to the case 𝑛 = 2. Let 𝑋, 𝑌 be
two sets and 𝑥, 𝑦 be two variables taking their values in 𝑋, 𝑌, respectively. A proposition
𝑅(𝑥, 𝑦) containing 𝑥, 𝑦 is called a relation or a binary relation if it can be determine
whether 𝑅(𝑎, 𝑏) is true or false for each pair (𝑎, 𝑏) in a Cartersian product 𝑋 × 𝑌. For
example, if both 𝑋 and 𝑌 are the set of rational integers, then the following propositions
relations: 𝑥 ≤ 𝑦, 𝑥 − 𝑦 is even, 𝑥 divides 𝑦. A relation 𝑅(𝑥, 𝑦) is sometimes written as
𝑥𝑅𝑦.

For a given relation 𝑅, we define its inverse relation 𝑅 −1 by 𝑦𝑅 −1 𝑥 ⟺ 𝑥𝑅𝑦. Then 𝑅 is

the inverse relation of 𝑅 −1 . In the example above, the inverse relation of 𝑥 ≤ 𝑦 is 𝑦 ≥ 𝑥,
and the inverse relation of 𝑥 is a divisor of 𝑦 is 𝑦 is a multiple of 𝑥. A relation 𝑅 is called
reflexive if 𝑥𝑅𝑥 holds. 𝑅 is called symmetric if 𝑥𝑅𝑦 ⟺ 𝑦𝑅𝑥 (namely, if 𝑅 and 𝑅 −1 are
identical). 𝑅 is called transitive if 𝑥𝑅𝑦 and 𝑦𝑅𝑧. 𝑅 is called antisymmetric if 𝑥𝑅𝑦 and 𝑦𝑅𝑥
imply 𝑥 = 𝑦. A reflexive, symmetric, and transitive relation is called an equivalence
relation. A reflexive and transitive relation is called a preordering. A reflexive, transitive,
and antisymmetric relation is called an ordering.

Suppose that we are given a relation 𝑥𝑅𝑦(𝑥 ∇ 𝑋, 𝑦 ∇ 𝑌). Then the set 𝐺 =
𝑥, 𝑦 𝑥𝑅𝑦 ,which consists of elements (𝑥, 𝑦) of the Cartesian product 𝑋 × 𝑌 satisfying
𝑥𝑅𝑦, is called the graph of the relation 𝑅.

RELATIONS BETWEEN 𝐏𝐧 𝐱 AND 𝐐𝐧 𝐱 :

(i) x 2 − 1 Qn Pn′ − Pn Qn ′ = c
Q n (x) ∞ dx
(ii) = ∫∝
P n (x) x 2 −1 Pn 2 x
1
(iii) Pn Qn−1 − Qn Pn−1 = n .
(2n−1)
(iv) Pn Qn−2 − Qn Pn−2 = n(n−1) x.

RELATIVE DISPLACEMENT: If during any interval of time P moves to P’ and Q moves to

Q’, then the relative displacement of P to Q during that interval is given by the change in
their relative position during that interval. It is determined by the vector difference
𝑃′ 𝑄 ′ −𝑃𝑄

RELATIVE SCALAR: A relative tensor of order zero is called relative scalar. Thus

𝜔
′
𝜕𝑥
𝐼 =𝐼
𝜕𝑥 ′
Then 𝐼 is called relative scalar of weight 𝜔. A relative scalar of weight one is called scalar
density. A relative scalar of weight zero called absolute scalar.

RELATIVE TENSOR: If the quantities 𝐴𝑖𝑗 satisfy the following transformation law

𝜔
𝜕𝑥 𝑎 𝜕𝑥 𝛽 𝜕𝑥
𝐴′𝑖𝑗 = 𝐴𝑎𝛽
𝜕𝑥 ′𝑖 𝜕𝑥 ′𝑗 𝜕𝑥 ′

Then 𝐴𝑖𝑗 is called a relative tensor of weight 𝜔. A relative tensor of weight one is called
tensor density. If the weight of a relative tension is zero, then the relative tensor is
called an absolute tensor.

RELATIVE VECTOR: A relative tensor of rank one is called relative vector. Thus if

𝜔
𝜕𝑥 𝑎 𝜕𝑥
𝐴′𝑖𝑗 = 𝐴𝑎𝛽
𝜕𝑥 ′𝑖 𝜕𝑥 ′

Then 𝐴𝑖 is called a relative vector of weight 𝜔. A relative vector of weight of one is celled
vector density. A relative vector of weight zero is called absolute vector.

RELATIVE VELOCITY: The relative velocity of Q with respect to P is the rate of change of
the position of Q relative to P and therefore, if uniform, is determined by the relative
displacement in one second.

REMAINDER THEOREM: Let 𝑘 be an integral domain, 𝑓 𝑋 𝜖𝑘[𝑋],and let 𝑔 𝑋 = 𝑋 −

𝛼(𝛼𝜖𝑘). Then using the division algorithm, we get

𝑓 𝑋 = 𝑋 − 𝛼 𝑞 𝑋 + 𝑟,

𝑞 𝑋 𝜖𝑘 𝑋 , 𝑟𝜖𝑘.

Therefore, 𝑓 𝛼 = 𝑟; that is, the remainder of 𝑓 𝑋 divided by 𝑋 − 𝛼 is equal to 𝑓 𝛼 .

This is called the remainder theorem. If 𝑓 𝛼 = 0, then 𝑓 𝑋 is divisible by 𝑋 − 𝛼 in
𝑘𝑋.

RENEGING: A customer may leave the queue due to impatience after joining it.

REPER: Reper is a collection of (maximally independent) unit vectors emanating from

the same end point. The number of vectors in a reper equals the dimension of the space.
Think of a cube. Remove all edges, except of some three that share a common vertex.
REPLACEMENT MODEL: There are two typical cases. One is the preventive maintenance
model, which is suitable when replacements are done under a routine policy because a
replacement or a repair before a failure is more effective than after a failure.
Probabilistic treatments we need to compare costs of both used and are mainly used,
and this model resembles those for queues and Markov processes. The other is a mode1
for deciding whether to replace a piece of equipment in use. In this case, types of
present cost are important.
REPLACEMENT PROBLEM: The problem of replacement is experienced in systems
where machines, individuals or capital assets are the main job performing main units.
The special feature of these units is that their level of performance / efficiency decrease
with time and one has to formulate some suitable replacement policy regarding these
units to keep the system up to some desired level of performance.

RESIDUE AT A POLE: Suppose a single valued function 𝑓(𝑧) has a pole of order 𝑚 at
𝑧 = 𝛼, then by definition of pole the principles part of Laurent’s expansion of 𝑓(𝑧)
contains only 𝑚 terms so that

∞ 𝑚
𝑛
𝑓 𝑧 = 𝑎𝑛 (𝑧 − 𝛼) + 𝑏𝑛 (𝑧 − 𝛼)𝑛
𝑛=0 𝑛 =1

αn = 2πi ∫ (z−α)n +1 , bn = 2πi ∫ (z−α)−n +1

1 f z dz 1 f z dz
where
C C

C being a circle z − α = r

1
Evidently b1 = 2πi ∫ c f z dz

The coefficient b1 is called the residues of 𝑓 z at the pole z = α and is denoted by the
symbol Res (z = α) = b1 .

RESIDUE AT INFINITY: If 𝑓 𝑧 has an isolated singularity at 𝑧 = ∞ or is analytic there,

their the residue at 𝑧 = ∞ is defined, as

1
Res (𝑧 = ∞) = − 2𝜋𝑖 ∫ 𝑐 𝑓 𝑧 𝑑𝑧

Where 𝐶 is any closed contour which encloses all the finite singularities of 𝑓(𝑧). The
integral is taken in positive direction (anticlockwise direction).
Cauchy’s Residue Theorem: If 𝑓(𝑧) is analytic within and on a closed contour 𝐶, except
at a finite number of poles z1 , z2 , z3 … , zn within 𝐶, then

c f z dz = 2πi Res (z = zr )
r=1

Where R.H.S. denotes sum of residues of 𝑓(𝑧) at its poles lying within 𝐶.

RESIDUE THEOREM: Residue theorem, also called Cauchy's residue theorem is a

powerful tool to evaluate line integrals of analytic functions over closed curves; it can
often be used to compute real integrals as well. It generalizes the Cauchy integral
theorem and Cauchy's integral formula.

Suppose 𝑈 is a simply connected open subset of the complex plane, and 𝑎1 , . . . , 𝑎𝑛 are
finitely many points of 𝑈 and 𝑓 is a function which is defined and holomorphic on 𝑈 \
{𝑎1 , . . . , 𝑎𝑛 }. If 𝛾 is a closed rectifiable curve in 𝑈 which does not meet any of the 𝑎𝑘 ,

If 𝛾 is a positively oriented simple closed curve, 𝐼(𝛾, 𝑎𝑘 ) = 1 if 𝑎𝑘 is in the interior of 𝛾,

and 0 if not, so

with the sum over those 𝑘 for which 𝑎𝑘 is inside 𝛾.

RESOLVENT EQUATION: Let 𝑓 𝜆 be the resolvent of the element 𝑓 𝜖 𝐴 (Banach

algebra). If 𝜆1 𝑎𝑛𝑑 𝜆2 are both in 𝜌 𝑓 , then
𝑓(𝜆1 )−1 𝑓 𝜆2 = 𝑓 − 𝜆1 𝑒 𝑓 𝜆2
= 𝑓 − 𝜆2 𝑒 𝑓 𝜆2 + 𝜆2 𝑒 − 𝜆1 𝑒 𝑓 𝜆2
= 𝑒 + 𝜆2 − 𝜆1 𝑓 𝜆2
𝑓 𝜆2 = 𝑓 𝜆1 + 𝜆2 − 𝜆1 𝑓 𝜆1 𝑓 𝜆2 .
Or
This implies
𝑓 𝜆1 − 𝑓 𝜆2 = 𝜆1 − 𝜆2 𝑓 𝜆1 𝑓 𝜆2
This relation is called the resolvent equation
RESOLVENT OF AN OPERATOR: The resolvent of an operator T is the operator valued
function defined on the resolvent set by the formula 𝑅(𝜆, 𝑇) = (𝑇 − 𝜆 𝐼)−1 .

Note that

 If | 𝜆 | > ||𝑇|| then 𝜆 ∇ 𝜌(𝑇), hence the spectrum is bounded.

 The resolvent set 𝜌(𝑇) is open, i.e for any 𝜆 ∇ 𝜌(𝑇) then there exist є > 0
such that all µ with | 𝜆 − µ | < є are also in 𝜌(𝑇), i.e. the resolvent set is open
and the spectrum is closed. Both statements together imply that the
spectrum is compact.
 𝑅(𝜆, 𝑇) − 𝑅(µ, 𝑇) = (𝜆 − µ)𝑅(𝜆, 𝑇)𝑅(µ, 𝑇) (First resolvent identity)
−1 −1
 𝑇−µ𝐼 → 𝑇−𝜆𝐼 𝑎𝑠 µ → 𝜆.
 For 𝑧 ∇ 𝜌(𝑡) the complex derivative 𝑑/𝑑𝑧 𝑅(𝑧, 𝑇) of the resolvent 𝑅(𝑧, 𝑇) is
well defined, i.e. the resolvent is an analytic function operator valued function
of 𝑧.
 The spectrum is non-empty.

RESOLVENT OPERATOR: Let A be a Banach algebra and 𝜌(𝑓) be the resolvent set of 𝑓.
The function
𝑓 ∶ 𝜌 𝑓 ⟶ 𝐴 𝑑𝑒𝑓𝑖𝑛𝑒𝑑 𝑏𝑦
𝑓 𝜆 = (𝑓 − 𝜆𝐼)−1
Is called resolvent operator of 𝑓.
The resolvent operator associated with 𝑓 is also called resolvent of 𝑓.
Note: 𝑓 𝜆 is a continuous function of 𝜆.
RESOLVENT SET: The resolvent set of 𝑓 is the complement of 𝜍 (𝑓) and it is denoted by
𝜌 𝑓 . 𝜌(𝑓) is an open subset of the complex plane which contains 𝑧: 𝑧 > 𝑓
RESOLVENT SET ρ T OF AN OPERATOR: The resolvent set 𝜌(𝑇) of an operator 𝑇 is the
set 𝜌 (𝑇) = {𝜆 ∇ ℂ: 𝑇 − 𝜆 𝐼 𝑖𝑠 𝑖𝑛𝑣𝑒𝑟𝑡𝑖𝑏𝑙𝑒}.

 The spectrum of operator 𝑇 ∇ 𝐵(𝐻), denoted 𝜍(𝑇), is the complement of the

resolvent set 𝜌(𝑇): 𝜍(𝑇) = {𝜆 ∇ ℂ: 𝑇 − 𝜆 𝐼 𝑖𝑠 𝑛𝑜𝑡 𝑖𝑛𝑣𝑒𝑟𝑡𝑖𝑏𝑙𝑒}.
 The spectrum 𝜍(𝑇) of a bounded operator 𝑇 is a nonempty compact (i.e. closed
and bounded) subset of ℂ.
 Let 𝐴 ∇ 𝐵(𝐻). If ||𝐴|| < 1 then 𝐼 − 𝐴 is invertible in 𝐵(𝐻) and inverse is given by
the Neumann series (𝐼 − 𝐴) − 1 = 𝐼 + 𝐴 + 𝐴2 + 𝐴3 + ⋯.
REVERSE CONVEX CONSTRAINT QUALIFICATION: Let X 0 be an open set in Rn and let 𝑔
be an 𝑚- dimensional vector function defined on X 0 .

Let 𝑋 = x: x ∇ X 0 , g(x) ≤ 0 . 𝑔 is said to satisfy the reverse convex constnat

qualification at x ∇ X if 𝑔 is differentiable at x, and if each i ∇ I either g i is linear of Rn ,
where

I = i gi x = 0

RICCI’S LEMMA: The metric tensors are covariant constants with respect to Christoffel
symbols.

RICCI’S THEOREM: The covariant derivatives of the tensors 𝑔𝑖𝑗 , 𝑔𝑖𝑗 , 𝑔𝑗𝑖 all vanish
identically.

RIEMANNIAN METRIC: A formula which expresses the distance between adjacent

points is called a line element or metric.

Examples : (i) 𝑑𝑠 2 = 𝑑𝑥 2 + 𝑑𝑦 2 .

(ii) 𝑑𝑠 2 = 𝑑𝑟 2 + 𝑟 2 𝜃 2 .

(iii) 𝑑𝑠 2 = 𝑑𝑥 2 + 𝑑𝑦 2 + 𝑑𝑧 2 .

This expresses the distance between (𝑥, 𝑦, 𝑧) and (𝑥 + 𝑑𝑥, 𝑦 + 𝑑𝑦, 𝑧 + 𝑑𝑧) in rectangular
Cartesian axes.

(iv) 𝑑𝑠 2 = 𝑎 𝑑𝑥 2 + 𝑏 𝑑𝑦 2 + 𝑐 𝑑𝑧 2 + 2𝑓 𝑑𝑥 𝑑𝑦 + 2𝑕 𝑑𝑧 𝑑𝑥.

This expresses the distance between curvilinear points (𝑥, 𝑦, 𝑧) and (𝑥, 𝑑𝑥, 𝑦 + 𝑑𝑦, 𝑧 +
𝑑𝑧) when the axes are not rectangular. Here the coefficients 𝑎, 𝑏, 𝑐, 𝑓, 𝑔, 𝑕 are functions
of co-ordinates 𝑥, 𝑦, 𝑧. Riemann extended this idea to a space of 𝑛-dimension and
defined the distance 𝑑𝑠 between adjacent points whose co-ordinates in any system are
𝑥 𝑖 and 𝑥 𝑖 + 𝑑𝑥 𝑖 by the formula, 𝑑𝑠 2 = 𝑔𝑖𝑗 𝑑𝑥 𝑖 𝑑𝑦 𝑗

Where g ij are functions of co-ordinates 𝑥 𝑖 .

The quadratic differential form on the R.H.S. of (1) is called Riemannian metric for n-
dimention al space. This differential dorm is assumed to be potistive definite. A space
characterized by this metric is called Riemannian space of n- dimentsions and is denoted
by ‘𝑉𝑛 ’. Geometry based on this metric is called Riemannian Geometry of 𝑛- dimensions.

RIEMANN INTEGRABLE FUNCTION: Let 𝑓 be a bounded function on a bounded interval

[𝑎, 𝑏]. 𝑓 is said to be Riemann integrable on [𝑎, 𝑏] if lower Riemann integral of 𝑓
coincides with upper Riemann integral of 𝑓 i.e.
𝑏 𝑏

𝑓𝑑𝑥 = 𝑓𝑑𝑥
𝑎 𝑎

and the common value of the integral is denoted as

𝑏

𝑓𝑑𝑥
𝑎

RIEMANN-LEBESGUE LEMMA: Consider the real Hilbert space of real valued square-
summarable function 𝐿20,2𝜋 . Then

𝟏 𝒔𝒊𝒏 𝒕 𝒄𝒐𝒔 𝒕 𝒔𝒊𝒏 𝟐𝒕 𝒄𝒐𝒔 𝟐𝒕

𝑨= ,
, , , ,…
𝟐𝝅 𝝅 𝝅 𝝅 𝝅

= 𝒇𝟎 , 𝒇𝟏 , 𝒇𝟐, …

1 𝑛
where 𝑓𝑛 = cos 𝑡 , 𝑛 even
𝜋 2

1 𝑛+1
= sin 𝑡 , 𝑛 odd
𝜋 2

In an orthonormal set in this space.

We define, for 𝑓 ∇ 𝐿2[0,2𝜋] , 𝑎𝑛 = 𝑓, 𝑓𝑛 ;

1 2𝜋
i.e. 𝑎0 = ∫0 𝑓 𝑡 𝑑𝑡,
(2𝜋)

1 2𝜋
𝑎1 = ∫0 𝑓 𝑡 𝑡 𝑑𝑡,
𝜋

⋯ ⋯ ⋯ ⋯ ⋯

If the series of the ordinary Fourier coefficients of 𝑓 be written, 𝑣𝑖𝑧.,

1 2𝜋
𝑎0 = ∫0 𝑓 𝑡 𝑑𝑡,
𝜋
 1 2𝜋 1
And for 𝑛 > 0, 𝑎𝑛 = ∫ 𝑓 𝑡 cos 𝑛𝑡 𝑑𝑡, 𝑛 even,
𝜋 0 2

1 2𝜋 1
and , 𝑏𝑛 = ∫ 𝑓 𝑡 sin 𝑛 + 1 𝑑𝑡, 𝑛 odd,
𝜋 0 2

then we find at once the following relationships:

𝛼0 = 𝑎0 𝜋/2, 𝛼1 = 𝑏1 𝜋, 𝛼2 𝜋, …

Since by Baseel’s inequality, we have

∞
2 2
𝛼𝑛 ≤ 𝑓 ,
𝑛=0

We see that

∞ 2𝜋
𝜋
𝛼02 + 𝜋 𝑎𝑛2 + 𝑏𝑛2 ≤ 𝑓2 𝑡 𝑑𝑡,
2 0
𝑛=0

Which implies

limn→∞ 𝛼𝑛 = 0 and limn⟶∞ 𝑏𝑛 = 0⁡

The above result is called the Riemann-Lebesgue Lemma.

RIEMANN-LEBESGUE THEOREM: A necessary and sufficient condition for a function 𝑓 to

be Riemann integrable on an interval [a, b] is that 𝑓 is bounded and that its set of points
of discontinuity in [a, b] forms a set of Lebesgue measure zero.
RIEMANN MAPPING THEOREM: If 𝑈 ≠ ∅, 𝐶 is a simply connected open subset of 𝐶 then
there is a biholomorphism 𝑓 ∶ 𝑈 → 𝐷 onto the open unit disc 𝐷 = {𝑧 ∇ 𝐶 ∶ |𝑧| < 1}.
RIEMANN REARRANGEMENT THEOREM: If an infinite series is conditionally convergent,
then its terms can be arranged in a permutation so that the new series converges to any
given value, or diverges.
RIEMANN REMOVABLE SINGULARITIES THEOREM: Let 𝑓 ∶ 𝐷(𝑃, 𝑟) \ {𝑃} → 𝐶 be
holomorphic and bounded. Then
(a) lim𝑧→𝑃 𝑓(𝑧) exists.
(b) The function 𝑓: 𝐷(𝑃, 𝑟) → 𝐶 defined by
( 𝑓(𝑧) 𝑖𝑓 𝑧 ≠ 𝑃
𝑓(𝑧) = is holomorphic.
lim𝜁 →𝑃 𝑓(𝜁) 𝑖𝑓 𝑧 = 𝑃
RIEMANN-ROCH THEOREM: For a divisor class 𝐷 on a closed Riemann surface ℛ of
genus 𝑔 and for an integer 𝑛, we have
𝑑𝑖𝑚(𝐷 + 𝑛𝑊) − 𝑑𝑖𝑚(−𝐷 − (1 − 𝑛)𝑊) = 𝑑𝑒𝑔𝐷 + (2𝑛 − 𝑙)(𝑔 − 𝑙),
where 𝑊 is the canonical divisor class.
RIEMANN-STIELTJES INTEGRABLE FUNCTION: Let 𝑓 be a bounded function on a
bounded interval [𝑎, 𝑏] and let 𝑔 be a monotonically non-decreasing function on [𝑎, 𝑏]. 𝑓
is said to be Riemann-Stieltjes integrable on [𝑎, 𝑏] if lower Riemann-Stieltjes integral of
𝑓 coincides with upper Riemann-Stieltjes integral of 𝑓 i.e.
𝑏 𝑏

𝑓𝑑𝑔 = 𝑓𝑑𝑔
𝑎 𝑎

RIEMANNIAN SUBMANIFOLD: If an immersion (or an embedding) 𝑓 of a Riemannian

manifold (𝑀, 𝑔) into a Riemannian manifold (𝑀, 𝑔) satisfies the condition 𝑓 ∗ 𝑔 = 𝑔,
then 𝑓 is called an isometric immersion (or embedding) and 𝑀 is called a Riemannian
submanifold of 𝑀.

𝑛 −𝜍 1 1 1
RIEMANN ZETA FUNCTION: The series 𝑛=1 𝑛 = 1𝜍 + 2𝜍 + 3𝜍 + … …

is uniformly convergent for all real 𝜍 ≥ 𝜍0 where is 𝜍0 is a fixed number such that
𝜍0 > 1. The series (1) is a majorant of the series.

𝜃 −𝑧
𝜁 𝑧 = 𝑛=1 𝑛 (where z = x + iy)

The function 𝜁(𝑧) is known as Riemann’s 𝜁 − 𝑓𝑢𝑛𝑐𝑡𝑖𝑜𝑛, where Re 𝑧 > 1. It plays a

central role in the application of complex analysis to number theory.

RIESZ–FRÉCHET LEMMA: Let H be a Hilbert space and α a continuous linear functional

on 𝐻, then there exists the unique 𝑦 ∇ 𝐻 such that 𝛼(𝑥) = ⟨ 𝑥, 𝑦 ⟩ for all 𝑥 ∇ 𝐻. Also
||𝛼||𝐻 ∗= ||𝑦||𝐻 .

RIESZ–FISHER THEOREM: Let ⌌𝑒𝑛 ⌍1∞ be an orthonormal sequence in a Hilbert space 𝐻.

∞ ∞
Then 𝑛=1 𝜆𝑛 𝑒𝑛 converges in 𝐻 if and only if 𝑛=1 | 𝜆𝑛 |2 < ∞.

RIESZ REPRESENTATION THEOREM: Let K be a compact Hausdorff space, and let

𝜆 ∇ 𝐶𝐾 (𝐾)∗ . There exists a unique µ ∇ 𝑀𝐾 (𝐾) such that λ(f) = ∫𝑋 𝑓𝑑𝜇. Furthermore,
||𝜆|| = ||µ||.
RIESZ REPRESENTATION THEOREM: Let 𝐻 be a Hilbert space, and let 𝐻 ∗ denote its dual
space, consisting of all continuous linear functionals from 𝐻 into the field 𝑹 or 𝑪. If 𝑥 is
an element of 𝐻, then the function 𝜑𝑥 , for all 𝑦 in 𝐻 defined by

where denotes the inner product of the Hilbert space, is an element of 𝐻 ∗ . The
Riesz representation theorem states that every element of 𝐻 ∗ can be written uniquely in
this form.

RIGHT CIRCULAR CONE OF SEMI-VERTICAL ANGLES 𝜶: Let the line OZ be taken as axis
of the cone and let 𝑃 = (𝑥, 𝑦, 𝑧) be any point on the surface of the cone. Further let 𝑢 be
the distance of 𝑃 from z-axis and let 𝑣 be the inclination of the plane containing 𝑃 and
the 𝑧 − axis to zx-plane. Then as shown in the above figure, we have

𝑥 = 𝑢 cos 𝑣, 𝑦 = sin 𝑣, 𝑧 = 𝑂𝐾 = 𝑢 cot 𝛼.

⟹ 𝑃𝑜𝑠𝑖𝑡𝑖𝑜𝑛 𝑣𝑒𝑐𝑡𝑜𝑟 𝑜𝑓 𝑃 = (𝑢 cos 𝑣, 𝑢 sin 𝑣, 𝑢 cot 𝛼

⟹ 𝑟 = (𝑢 cos 𝑣, 𝑢 sin 𝑣, 𝑢 𝑐𝑜𝑡 𝛼

⟹ 𝑟1 = (cos 𝑣, sin 𝑣, cot 𝛼); 𝑟2 = (−𝑢 sin 𝑣, 𝑢 cos 𝑣, 0)

⟹ 𝑟1 . 𝑟2 = 0 ⟹ 𝑃𝑎𝑟𝑎𝑚𝑒𝑡𝑟𝑖𝑐 𝑐𝑢𝑟𝑣𝑒𝑠 𝑎𝑟𝑒 𝑜𝑟𝑡𝑕𝑜𝑔𝑜𝑛𝑎𝑙

𝑖 𝑗 𝑘
Here 𝑟1 × 𝑟2 = cos 𝑣 sin 𝑣 cot 𝛼
−𝑢 𝑠𝑖𝑛 𝑣 𝑢 cos 𝑣 0

= (−𝑢 cos 𝑣 cot 𝛼) 𝑖 − 𝑗 (𝑢 sin 𝑣 cot 𝛼) + 𝑘(𝑢)

⟹ (𝑟1 . 𝑟2 )𝑢=0 = 0 𝑖. 𝑒., vertex is the only singularity.

RIGHT HELOCOID: The surface generated by the helicoids motion of a straight line
meeting the axis in perpendicular direction is called right-helicoid.

RIGHT NULLITY OF A MATRIX: Suppose 𝑌 is an 𝑛-vector written in the form of a column

vector. Then the matrix product 𝐴𝑌 is defined. The subspaces 𝑇 or 𝑉𝑛 generated by the
column vector 𝑌 belonging to 𝑉𝑛 such that 𝐴𝑌 = 0 is called the column null space of the
matrix 𝐴. The dimension 𝑡 of 𝑇 is called the right nullity or column nullity of the matrix
𝐴.
RING HOMOMORPHISM: Let 𝑅 and 𝑆 be rings. A ring homomorphism is a function
𝑓 ∶ 𝑅 → 𝑆 such that:
• 𝑓(𝑎 + 𝑏) = 𝑓(𝑎) + 𝑓(𝑏) for all 𝑎, 𝑏 ∇ 𝑅
• 𝑓(𝑎 · 𝑏) = 𝑓(𝑎) · 𝑓(𝑏) for all 𝑎, 𝑏 ∇ 𝑅

RING ISOMORPHISM: A ring isomorphism 𝜙: 𝑅 → 𝑆 between rings 𝑅 and 𝑆 is a

homomorphism that is also a bijection between 𝑅 and 𝑆. The inverse of an isomorphism
is itself an isomorphism. Two rings are said to be isomorphic if there is an isomorphism
between them.

RINGS OF POLYNOMIALS: Let 𝑅 be a ring and let 𝑋1 , … , 𝑋𝑛 be variables (letters,

indetermionates, or symbols). Then the set of polynomials in 𝑋1 , … , 𝑋𝑛 with coefficients
in 𝑅 is called the ring of polynomials (or polynomial ring) in 𝑛 variables 𝑋1 , … , 𝑋𝑛 over 𝑅
and is denoted by 𝑅 𝑋1 , … , 𝑋𝑛 . On the other hand, when 𝑅 and 𝑅′ are rings with
common unity element and 𝑅 ⊂ 𝑅′, then for a subset 𝑆 of 𝑅′ we denote the subring of 𝑅′,
generated by 𝑆 over 𝑅, by 𝑅[𝑆]. When 𝑆 = 𝑥1 , … , 𝑥𝑛 , then there is a homomorphism 𝜑
of the polynomial ring 𝑅[𝑥1 , … , 𝑥𝑛 ] onto 𝑅[𝑆] defined by

𝑖 𝑖 𝑖 𝑖
𝜑 𝑎𝑖𝑗 , … , 𝑖𝑛 𝑋11 … 𝑋𝑛𝑛 = 𝑎𝑖𝑗 , … , 𝑖𝑛 𝑋11 … 𝑋𝑛𝑛

(𝑎𝑖𝑗 , … , 𝑖𝑛 ∇ 𝑅)

If 𝜑 is an isomorphism, then 𝑥1 , … , 𝑥𝑛 are said to be algebraically independent over 𝑅;

and otherwise, algebraically dependent over 𝑅.

RINGS: A nonempty set 𝐴 is called a ring if the following conditions are satisfied.

(1) Two operations, called addition and multiplication (the ring operations), are
defined, which send an arbitrary pair of elements 𝑎, 𝑏 of 𝐴 to elements 𝑎 + 𝑏 and 𝑎𝑏
of 𝐴.
(2) For arbitrary elements 𝑎, 𝑏, 𝑐 of 𝐴, these operations satisfy the following four laws;
(i) 𝑎 + 𝑏 = 𝑏 + 𝑎 (commutative law of additions); (ii) 𝑎+𝑏 +𝑐 =𝑎+
𝑏 + 𝑐 , 𝑎𝑏 𝑐 = 𝑎(𝑏𝑐)(associative laws); (iii) 𝑎 𝑏 + 𝑐 = 𝑎𝑏 + 𝑎𝑐, 𝑎 + 𝑏 𝑐 = 𝑎𝑐 +
𝑏𝑐 (distributive laws); and (iv) for every pair 𝑎, 𝑏 of elements of 𝐴, there exists a
unique element 𝑐 of 𝐴 such that 𝑎 + 𝑐 = 𝑏.
Thus a ring 𝐴 is an A belian group under addition. Each element 𝑎 of a ring 𝐴
determines operations 𝐿𝑎 and 𝑅𝑎 defined by 𝐿𝑎 𝑥 = 𝑎𝑥, 𝑅𝑎 𝑥 − 𝑥𝑎(𝑥 ∇ 𝐴). Thus a
ring 𝐴 has the structure of a left 𝐴- module and a right 𝐴- module. Since the operations
𝐿𝑎 and 𝑅𝑏 commute foe every pair 𝑎, 𝑏 of elements of 𝐴, the ring 𝐴 is also an 𝐴 − 𝐴-
bimodule.

The identity elements of 𝐴 under addition is called the zero element and is denoted by 0.
It satisfies the equation 𝑎0 = 0𝑎 = 0(𝑎 ∇ 𝐴). An element 𝑒 ∇ 𝐴 is called a unity element
(identity element or unit element) of 𝐴 if it satisfies 𝑎𝑒 = 𝑒𝑎 = 𝑎(𝑎 ∇ 𝐴). If 𝐴 has such a
unity element, it is unique and is often denoted by 1. A ring with unity element is called
a unitary ring. Most of the important examples of rings are unitary.

Hence we often call a unitary ring simply a ring. If a ring has only one member (namely,
0) then 0 is the unity element of the ring. Such a ring is called a zero ring. However, if a
ring has more than one element, the unity element is distinct from the zero elements. A
ring is called a commutative ring if it satisfies 𝑣 𝑎𝑏 = 𝑏𝑎 (𝑎, 𝑏 ∇ 𝐴)(commutative law
for multiplication.

An element 𝑎 ≠ 0 of a ring 𝐴 is called a zero divisor if there exists an element 𝑏 ≠ 0 such

that 𝑎𝑏 = 0 or 𝑏𝑎 = 0. A commutative unitary ring having more than one element is
called an integral domain if it has no zero divisors Elements 𝑎 and 𝑏 of a ring are said to
be orthogonal if 𝑎𝑏 = 𝑏𝑎 = 0. An element 𝑎 satisfying 𝑎𝑛 = 0 for some positive integer
𝑛 is called a nilpotent element, and a nonzero element 𝑎 satisfying 𝑎2 = 𝑎 is called an
idempotent element. An idempotent element is said to be primitive if it cannot be
represented as the sum of two orthogonal idempotent elements. For any subsets 𝑆 and
𝑇 of a ring 𝐴, let 𝑆 + 𝑇(𝑆𝑇) denote the set of elements 𝑠 + 𝑡 𝑠𝑡 (𝑠 ∇ 𝑆, 𝑡 ∇ 𝑇). In
particular, 𝑆𝑆 is denoted by 𝑆 2 (similarly for 𝑆 3 , 𝑆 4 etc)., and furthermore, 𝑎 + 𝑆 𝑎 𝑆)is
denoted by 𝑎 + 𝑆(𝑎𝑆). If 𝑆𝑇 = 𝑇𝑆 = 0 , then subsets 𝑆 and 𝑇 are said to be orthogonal.
A subset 𝑆 of a ring is said to be nilpotent if 𝑆 𝑛 = 0 for some positive integer 𝑛, and
idempotent if 𝑆 2 = 𝑆.

Homomorphism: A mapping 𝑓: 𝐴 → 𝐵 of a ring 𝐴 into a ring 𝐵 satisfying conditions (i)

𝑓 𝑎 + 𝑏 = 𝑓 𝑎 + 𝑓(𝑏) and (ii) 𝑓 𝑎𝑏 = 𝑓 𝑎 𝑓 𝑏 (𝑎, 𝑏 ∇ 𝐴) is called a homomorphism.
If a homomorphism. If a homomorphism 𝑓 is bijective, then the inverse mapping
𝑓 −1 : 𝐵 → 𝐴 is also a homomorphism, and in this case 𝑓 is called an isomorphism. More
precisely, a homomorphis, (isomorphism) of rings is often called a ring homomorphism
(ring isomorphism). There exists only one homomorphism of any ring onto the zero
ring. For unitary rings 𝐴 and 𝐵, a homomorphism 𝑓: 𝐴 → 𝐵 is said to be unitary if it
maps the unity element of 𝐴 to the unity element of 𝐵. By a homomorphism, a unitary
homomorphism is usually meant. In this sense there exists a unique homomorphism of
ring 𝑍 of rational integers into an arbitrary unitary ring. A composite of
homomorphisms is also a homomorphism. The identity mapping 𝑙𝐴 of a ring 𝐴 is an
isomorphism. A homomorphism of a ring 𝐴 into itself is called an endomorphism, and
an isomorphism of 𝐴 onto itself is called an automorphism of 𝐴. If 𝑎 is an invertible
element of a unitary ring 𝐴, then the mapping 𝑥 − 𝑎𝑥𝑎−1 (𝑥 ∇ 𝐴) is an automorphism of
𝐴, called an inner atuomorphism.

RODRIGUE’S FORMULA DIFFERENTIAL GEOMETRY): A necessary and sufficient

condition that a curve on a surface be line of curvature is that

𝑑𝑁 𝑑𝑟
∝ 𝑜𝑟 𝑑𝑁, 𝜅𝑑𝑟 = 0
𝑑𝑠 𝑑𝑠

at each of its points, where 𝜅 denotes the normal curvature.

RODRIGUES FORMULAS FOR LAGUERRE POLYNOMIALS: Rodrigues formulas for

Laguerre polynomials is

𝑒 𝑥 𝑑𝑛
𝐿𝑛 𝑥 = 𝑥 𝑛 𝑒 −𝑥
𝑛 𝑑𝑥 𝑛

ROLLE’S THEOREM: If a real valued function 𝑓(𝑥) is such that

 𝑓(𝑥) is continuous in the closed interval [𝑎, 𝑏]
 𝑓(𝑥) is differentiable in the open interval (𝑎, 𝑏)
 𝑓(𝑎) = 𝑓(𝑏)
then there exist at least one value of 𝑥 = 𝑐 ∇ (𝑎, 𝑏) such that 𝑓 ′ (𝑐) = 0.
ROMAN NUMERALS: Roman numerals are a method of writing numbers employed
primarily by the ancient Romans. It place of digits, the Romans used letters to represent
the numbers central to the system:
I ----------1
V----------5
X---------- 10
L---------- 50
C---------- 100
D---------- 500
M---------- 1000
Larger numbers can be made by writing a bar over the letter, which means one
thousand times as much. For instance 𝑉 is 5000.
ROOTS OF UNITY: A root of unity is a solution of 𝑧 𝑛 = 1, with 𝑧 ∇ 𝐶 and 𝑛 a positive
integer.

ROTATION (COMPLEX ANALYSIS): By the transformation 𝑤 = 𝑧𝑒 𝑖𝜃0, figures in

𝑧 − 𝑝𝑙𝑎𝑛𝑒 are rotated through an angle 𝜃0 < 0, the rotation is clock wise.

ROUCHÉ'S THEOREM: If the complex-valued functions 𝑓 and 𝑔 are holomorphic inside

and on some closed contour 𝐾, with |𝑔(𝑧)| < |𝑓(𝑧)| on 𝐾, then 𝑓 and 𝑓 + 𝑔 have the
same number of zeros inside 𝐾, where each zero is counted as many times as
its multiplicity. This theorem assumes that the contour 𝐾 is simple, that is, without self-
intersections.
ROUNDING ERROR: Rounding error is the difference between a rounded-off numerical
value and the actual value. A rounded quantity is represented by a numeral with a fixed
number of allowed digits, with the last digit set to the value that produces the smallest
difference between the rounded quantity and the actual quantity.

Rounding can produce a value that is easier to deal with than the actual value, especially
if the actual value contains a lot of digits. Rounding can also be done to indicate the
relative precision of a value. For example, the irrational number pi equals
approximately 3.14, rounded to two decimal places or three significant digits.

As an example of rounding error, consider the speed of light in a vacuum. The official
value is 299,792,458 meters per second. In scientific (power-of-10) notation, that
quantity is expressed as 2.99792458 x 108. Rounding it to three decimal places yields
2.998 x 108. The rounding error is the difference between the actual value and the
rounded value, in this case (2.998 - 2.99792458) x 108, which works out to 0.00007542
x 108. Expressed in the correct scientific notation format, that value is 7.542 x 103,
which equals 7542 in plain decimal notation.
ROW AND COLUMN EQUIVALENCE OF MATRICES: A matrix 𝐴 is said to be row
equivalent to 𝐵 if 𝐵 is obtainable from 𝐴 by a finite number of 𝐸- row transformations of
R
𝐴. Symbolically, we then write A~B . Similarly a matrix 𝐴 is said to be column equivalent
to 𝐵 if 𝐵 is obtainable from 𝐴 by a finite number of 𝐸 -column transformation of 𝐴.
C
Symbolically, we then write A~B .

ROW MATRICES: Column matrices: Any 1 × n matric which has only one row and n
columns is called a row matrix or a row vector. Similarly any 𝑚 × 1 matrix which has 𝑚
rows and one column is a column matrix or a column vector.

For example, 𝑋 = 2 7 −8 5 11 1×5 is a row matrix of the type 1 × 5 while

2
Y= −9 is a column matrix of the type 3 × 1.
11 3×1

ROW RANK OF A MATRIX: Let A = aij be any 𝑚 × 𝑛 matrix. Each of the 𝑚 × 𝑛 matrix.
Each of the 𝑚 rows of 𝐴 consists of 𝑛 elements. Therefore the row vectors of 𝐴 are 𝑛-
vectors. These row vectors of 𝐴 will span a subspace 𝑅 to 𝑉𝑛 . This subspace 𝑅 is called
the row space of the matrix 𝐴. The dimension 𝑟 of 𝑅 is called the rows rank of 𝐴. The
dimension 𝑟 of 𝑅 is called the row rank of 𝐴. In other words the row rank of a matrix 𝐴 is
equal to the maximum number of linearly independent rows of 𝐴.

RSA CRYPTOGRAPHIC THEOREM: Let 𝑝 and 𝑞 be distinct prime numbers, let 𝑚 = 𝑝𝑞

and let 𝑠 = (𝑝 − 1)(𝑞 − 1). Let 𝑗 and 𝑘 be positive integers with the property that
𝑗 ≡ 𝑘 (𝑚𝑜𝑑 𝑠). Then 𝑥 𝑗 ≡ 𝑥 𝑘 (𝑚𝑜𝑑 𝑚) for all integers 𝑥.
RULED SURFACE (DEVELOPABLE AND SKEW): A surface which is generated by the
motion of the motion of one parameter family of straight lines is called a ruled surface
and the line is called its generating line or generators or ruling. Some special forms of
ruled surfaces are cones, cylinder and conicoids. Ruled surface are divided in two
distinct classes, one class of ruled surface is those on which consecutive generators
intersect and another class of ruled surface is those on which consecutive generators do
not intersect; these class of ruled surface are called developable and skew surface
respectively. Skew surface are also called scrolls.
RUSSEL’S APPROXIMATION METHOD: For each source row i remaining under
consideration, determine its 𝑢𝑖 , the largest unit cost 𝑐𝑖𝑗 still remaining in that row. For
each destination column 𝑗 remaining under consideration, determine its 𝑣𝑗 , which is the
largest unit cost 𝑐𝑖𝑗 still remaining in that column. For each variable 𝑥𝑖𝑗 not previously

selected in these rows and columns, calculate ∆𝑖𝑗 = 𝑐𝑖𝑗 – 𝑢𝑖 – 𝑣𝑗 . Select the variable
having the largest negative value of ∆𝑖𝑗 .

RUSSELL’S PARADOX: Suppose that for any coherent proposition 𝑃(𝑥), we can construct
a set {𝑥 ∶ 𝑃(𝑥)}. Let 𝑆 = {𝑥 ∶ 𝑥 ∈ 𝑥}. Suppose 𝑆 ∇ 𝑆; then, by definition, 𝑆 ∈ 𝑆.
Likewise, if 𝑆 ∈ 𝑆, then by definition 𝑆 ∇ 𝑆. Therefore, we have a contradiction.
Bertrand Russell gave this paradox as an example of how a purely intuitive set theory
can be inconsistent. The regularity axiom, one of the Zermelo-Fraenkel axioms, was
devised to avoid this paradox by prohibiting self-swallowing sets.
SADDLE POITNT: A point which is minimum in its row and maximum in its column is
known as the saddle point.

SAINT – VENANT’S PRINCIPLE: If some distribution of forces acing on the portion of the
surface of a body is replaced by a different distribution of forces acting on the same
portion of the body, then the effects of the two different distributions on the parts of the
body sufficiently far removed from the region of applications of the forces are
essentially the same, provided that the two distributions of forces are statically
equivalent. This principle is also called the principle of softening of boundary
conditions.

SCALAR: A scalar is a quantity that is invariant under coordinate transformation, also

known as a tensor of rank 0. For example, the number 1 is a scalar, so is any number or
variable 𝑛 ∇ 𝑅. The point (3, 4) is not a scalar because it is variable under rotation. As
such, a scalar can be an element of a field over which a vector space is defined.
SCALAR MATRIX: A diagonal matrix whose diagonal elements are all equal is called a
scalar matrix.

k 0 … 0
0 k … ⋮
If 𝑆 = 0 0 … ⋮
⋮ ⋮ … ⋮
0 0 … k

is an n-rowed scalar matrix each of whose diagonal elements equal to 𝑘 and 𝐴 is any n-
rowed square matrix, then
 𝐴𝑆 = 𝑆𝐴 = 𝑘𝐴

i.e., the pre-multiplication or the post- multiplication of 𝐴 by 𝑆 has the same effect as the
multiplication of 𝐴 by the scalar 𝑘. This is perhaps the motivation behind the name
‘scalar matrix’.

As a particular case, if we take

a11 a12 a13 k 0 0

A = a21 a22 a23 and S = 0 k 0 , then
a31 a32 a33 0 0 k

a11 a12 a13 k 0 0 ka11 ka12 ka13

a
A = 21 a22 a23 × 0 k 0 = ka21 ka22 ka23 = 𝑘𝐴
a31 a32 a33 0 0 k ka31 ka32 ka33

Similarly 𝑆𝐴 = 𝑘𝐴. Hence 𝑆𝐴 = 𝐴𝑆 = 𝑘𝐴.

SCALAR MULTIPLICATION OF LINEAR MAPS: We define a map

𝛼𝑇: 𝑈 → 𝑉

by the rule

(𝛼𝑇) 𝑢 = 𝛼 𝑇 𝑢 ∀ 𝑢 𝜖 𝑈

SCALAR PRODUCT OF FOUR VECTORS: If 𝒂, 𝒃, 𝒄, 𝒅 are four vectors, the products

𝒂 × 𝒃 . 𝒄 × 𝒅 , 𝒂 × 𝒅 . 𝒃 × 𝒄 etc. are called the scalar products of four vectors.

Note that

𝒂∙𝒄 𝒃∙𝒄
𝒂 ×𝒃 . 𝒄×𝒅 =
𝒂∙𝒅 𝒃∙𝒅

This relation is known as Lagrange’s Identity.

SCALAR TRIPLE PRODUCT: The scalar product of two vectors one of which is itself the
vector product of two vectors is a scalar quantity called a “Scalar Triple Product”. This if
𝒂, 𝒃 and 𝒄 be three vectors, then 𝒂 × 𝒃 . 𝒄 is called the scalar triple product of these
three vectors. Since the scalar triple product involves both the signs of ‘cross’ and ‘dot’,
therefore it is sometimes also called the mixed product.
SCATTER PLOT: A scatter plot is a set of points plotted on a horizontal and vertical axes.
Scatter plots are important in statistics because they can show the extent of correlation,
if any, between the values of observed quantities or phenomena (called variables). If no
correlation exists between the variables, the points appear randomly scattered on the
coordinate plane. If a large correlation exists, the points concentrate near a straight line.
Scatter plots are useful data visualization tools for illustrating a trend.

Besides showing the extent of correlation, a scatter plot shows the sense of the
correlation:

 If the vertical (or y-axis) variable increases as the horizontal (or x-axis) variable
increases, the correlation is positive.
 If the y-axis variable decreases as the x-axis variable increases or vice-versa, the
correlation is negative.
 If it is impossible to establish either of the above criteria, then the correlation is
zero.
The maximum possible positive correlation is +1 or +100%, when all the points in a
scatter plot lie exactly along a straight line with a positive slope. The maximum possible
negative correlation is -1 or -100%, in which case all the points lie exactly along a
straight line with a negative slope.

Correlation is often confused with causation, either accidentally (as a result of false or
unproved hypotheses) or deliberately (with intent to deceive). However, in the pure
sense, while a scatter plot can reveal the nature and extent of correlation, it says nothing
about causation.

SCHEDULING MODE1 (SCHEDULING AND PRODUCTION PLANNING): Network

scheduling is used to schedule complicated projects (for example, construction of
buildings) that consist of a large number of jobs related to each other in some natural
order. PERT (program evaluation and review technique) and CPM (critical path
method) are popular computational methods for this model. Job shop scheduling is used
when we have m jobs and n machines and each job requires some of the machines in a
given order. The mode1 allows us to fïnd an optimal order (in some certain sense) of
jobs to be implemented on each machine.
SCHOTTKY’S THEOREM: Let 𝛼 and 𝛽 be positive real numbers with 𝛽 < 1. Then there is
a constant 𝐶 = 𝐶𝛼 , 𝛽 > 0 with the following property. If 𝑓 ∶ ∆ = 𝐵(0, 1) → 𝐶\{0, 1} is
analytic, with |𝑓(0)| ≤ 𝛼, then |𝑓(𝑧)| ≤ 𝐶 for |𝑧| ≤ 𝛽.

SCHRÖDER–BERNSTEIN THEOREM: If there exist injective functions 𝑓: 𝐴 →

𝐵 and 𝑔 : 𝐵 → 𝐴 between the sets 𝐴 and 𝐵, then there exists a
bijective function 𝑕 : 𝐴 → 𝐵. In terms of the cardinality of the two sets, this means that
if |𝐴| ≤ |𝐵| and |𝐵| ≤ |𝐴|, then |𝐴| = |𝐵|; that is, 𝐴 and 𝐵 are equipollent. This is a
useful feature in the ordering of cardinal numbers. The theorem is also known as the
Cantor–Bernstein theorem, or the Cantor–Schroeder–Bernstein theorem.

SCHWARZ- CHRISTOFFEL THEOREM: The transformation from the 𝜉- plane to the 𝑧-

plane is defined by

𝑑𝑧 𝜃1 𝜃2 𝜃𝑛
= 𝐾(𝜁 − 𝑎1 ) 𝜋 −1 (𝜁 − 𝑎2 ) 𝜋 −1 … … … (𝜉 − 𝑎𝑛 ) 𝜋 −1
𝑑𝜉

Which transforms the real axis 𝜂 = 0 in the 𝜁(= 𝜉 + 𝑖𝜂) -plane into the boundary of a
closed polygon in the 𝑧(= 𝑥 + 𝑖𝑦)- plane in such a way that the vertices of the simple
closed polygon correspond to the points 𝑎1 , 𝑎2 , … , 𝑎𝑛 and the interior angles of the
polygon 𝜃1 , 𝜃2 , … , 𝜃𝑛 , 𝐾 is a constant which may be complex.

SCHWARZ LEMMA (COMPLEX ANALYSIS): Suppose

(i) 𝑓 𝑧 is analytic in domain defined by 𝑧 < 1

(ii) 𝑓(𝑧) ≤ 1, (iii) 𝑓 0 = 0

Then 𝑓 𝑧 ≤ 𝑧 and 𝑓 ′ (𝑧) < 1

Equality occurs only if 𝑓(𝑧) is a linear map 𝑤 = 𝑓 𝑧 = 𝑧𝑒 𝑖𝛼 , where is 𝛼 is a real

constant.

SCHWARZ REFLECTION PRINCIPLE: Suppose that 𝑓1 (𝑧) is analytic in the region 𝑅1 and
that 𝑓1 (𝑧) takes only real values on the part 𝐿𝑀𝑁 of the real axis. Then Schwarz’s
reflection principle states that the analytic continuation 𝑓2 (𝑧) of 𝑓1 (𝑧) into the domain
𝑅2 (considered as the mirror image of 𝑅1 with 𝐿𝑀𝑁 as mirror is given by 𝑓1 𝑧 .
SCREW- CURVATURE: The arc rate at which principal normal changes direction
dn
i. e. , ds is called the screw curvature vector and its magnitude is given by k 2 + τ2 .

SECOND-COUNTABLE SPACES: (X, T) is second countable if it has a countable base. Note

that every second countable space is first countable, Lindelöf as well as paracompact
space. The property is stable under taking a subspace, Cartesian product or countable
union. Rn, with open sets the open balls of rational radius and center, or all rational
rectangles are examples of second countable space.

SECOND ISOMORPHISM THEOREM: Let 𝑀 and 𝑁 be normal subgroups of a group 𝐺,

𝐺 𝐺/𝑀
where 𝑀 ⊂ 𝑁. Then 𝑁 ≅ .
𝑁/𝑀

SELF DUAL CODE: Let 𝐶 be a linear code. Then 𝐶 is self dual code if 𝐶 = 𝐶 ⊥

SELF ORTHOGONAL CODE: Let 𝐶 be a linear code. Then 𝐶 is self orthogonal code if
𝐶 ⊆ 𝐶⊥
SEMIGROUPS: A semigroup consists of a set on which is defined an associative binary
operation. We may denote by 𝐴,∗ a semigroup consisting of a set 𝐴 together with an
associative binary operation ∗ on 𝐴. A semigroup 𝐴,∗ is said to be commutative (or
Abelian) if the binary operation ∗ is commutative. The set of natural numbers, with the
operation of addition, is a commutative semigroup, as is the set of natural numbers with
the operation of multiplication. Let 𝐴,∗ be a semigroup. Given any element 𝑎 of 𝐴, we
define
𝑎1 = 𝑎,
𝑎2 = 𝑎 ∗ 𝑎,
𝑎3 = 𝑎 ∗ 𝑎2 = 𝑎 ∗ 𝑎 ∗ 𝑎 ,
𝑎4 = 𝑎 ∗ 𝑎3 = 𝑎 ∗ 𝑎 ∗ 𝑎 ∗ 𝑎 ,
𝑎 5 = 𝑎 ∗ 𝑎4 = 𝑎 ∗ (𝑎 ∗ (𝑎 ∗ (𝑎 ∗ 𝑎))), . ..
In general we define a n recursively for all natural numbers n so that a 1 = a and a n = a
∗ a n−1 whenever n > 1.
SEMIRING: A semiring S of sets is the collection such that

2. it is closed under intersection;

3. for 𝐴, 𝐵 ∇ 𝑆 we have 𝐴\𝐵 = 𝐶1 ∪ … ∪ 𝐶𝑁 with 𝐶𝑘 ∇ 𝑆.

Note that the following are semirings but not rings:

  The collection of intervals [𝑎, 𝑏) on the real line;

 The collection of all rectangles { 𝑎 ≤ 𝑥 < 𝑏, 𝑐 ≤ 𝑦 < 𝑑 } on the plane.

 Let S be a semiring. Then the collection of all disjoint unions ⋃𝑛𝑘=1 𝐴𝑘 , where
𝐴𝑘 ∇ 𝑆, is a ring. We call it, the ring 𝑅(𝑆) generated by the semiring 𝑆.

 Let S be a semiring. Then Any ring containing 𝑆 contains 𝑅(𝑆) is also a semiring.

SEMISIMPLE GROUPS AND REDUCTIVE GROUPS: In an algebraic group 𝐺 defined over 𝑘,

there exists a largest connected solvable closed normal subgroup 𝑅 , called the radical of
𝐺. The unipotent part 𝑅1 , of 𝑅 is called the unipotent radical of G. When 𝑅 = {𝑒}, 𝐺 is
called semisimple. When 𝑅 is a torus, namely, 𝑅1 = {𝑒}, 𝐺 is called reductive.
Semisimplicity and reductiveness are preserved under forming a direct product and
taking the image (or inverse image) of an isogeny.
SENSITIVITY ANALYSIS: Sensitivity Analysis deals with the effect on the optimal
solution of making changes in the values of the model parameters 𝑎𝑖𝑗 , 𝑏𝑖 𝑐𝑗 .

Procedure for sensitivity analysis (Outline of system):

 Revision of model
 Revision of final tableau
 Conversion to proper form
 Feasibility test
 Optimality test
 Re-optimize
Many network optimization models are actually special types of linear programming
problems.
Four important kinds of network problems are:
 the shortest path problem
 the minimum spanning tree problem
 the maximum flow problem
 the minimum cost flow problem
SEPARABLE FIELD EXTENSION: An algebraic field extension 𝐿: 𝐾 is said to be
separable over 𝐾 if the minimum polynomial of each element of 𝐿 is separable over 𝐾.
SEPARABLE POLYNOMIAL: Let 𝐾 be a field. An irreducible polynomial in 𝐾[𝑥] is said to
be separable over 𝐾 if it does not have repeated roots in a splitting field. A polynomial in
𝐾[𝑥] is said to separable over 𝐾 if all its irreducible factors are separable over 𝐾. A
polynomial is said to be inseparable if it is not separable.

SEQUENCING PROBLEM: We will consider the problem of performing 𝑛 jobs on each of

𝑚 machines. We are given the order of the machines for each job, in which it should go
to the machines. We are given the order of the machines for each job, in which it should
go to the machines. We also know the actual or expected time required by the jobs on
𝑚
each of the machines. Our problem is to find that sequence out of 𝑛! sequences
which minimizes the total elapsed time .i.e. the time from start of the job upto the
completion of the last job.

Mathematically, if we use the notations

𝐴𝑖 = time estimated for the 𝑖 𝑡𝑕 job on machine 𝐴, 𝑖 = 1,2 , … . 𝑛.

(Similarly we can interpret 𝐵𝑖 𝑎𝑛𝑑 𝐶𝑖 etc.)

T= the total elapsed time

Then we determine a sequence of jobs, i.e. a permutation of numbers 1, 2…,n for each
machine, which minimizes the time T.

All types of sequencing problems cannot be solved. The satisfactory solutions are
available only in few cases.

SEQUENTIAL SPACES: The space (𝑋, 𝑇) is sequential if for every open set 𝐴 ⊂ 𝑋 every
sequence convergent to a point in 𝐴 is eventually in 𝐴.
SERIES OF POSITIVE TERMS: Suppose that 𝑎𝑛 is a series of positive (or nonnegative)
terms. Since its partial sums 𝑠𝑛 form a ↑monotonically increasing sequence, the series is
∞ −𝑝
convergent if and only if 𝑠𝑛 is bounded. For example, the series 𝑛=1 𝑛 (𝑝 > 0)
converges if 𝑝 < 1 because 𝑠𝑛 < 2𝑝−1 (2𝑝−1 − 1), whereas it diverges if 𝑝 ≤ 1 because
∞ 𝑛−1
𝑠2𝑚 +1 > 1 + (𝑚 + 1) 2. The geometric series 𝑛=1 𝑎 (𝑎 > 0) converges for 𝑎 < 1
because 𝑠𝑛 = (1 − 𝑎𝑛 ) (1 − 𝑎), and diverges for 𝑎 ≥ 1 because 𝑠𝑛 ≥ 𝑛.

SERRET- FRENET FORMULAE: The following set of three relations involving space
derivations of fundamental unit vectors 𝑡, 𝑛, 𝑏 are known as Serret- Frenet Formulae.
 dt
= kn
ds

dn
= τb = kt
ds

db
= τn
ds
SERVICING TIME: It is the time taken for servicing a particular arrival.

SET FUNCTIONS: A function whose domain is a family of sets is called a set function.
Usually we consider set functions that take real values or ±∞. For example, if 𝑓(𝑥) is a
real-valued function defined on a set 𝑋, and if we assign to each subset 𝐴 of 𝑋 values
such as 𝑠𝑢𝑝𝐴 𝑓, 𝑖𝑛𝑓𝐴 𝑓, or 𝑠𝑢𝑝𝐴 𝑓 − 𝑖𝑛𝑓𝐴 𝑓, then we obtain a corresponding set function. In
particular, a set function whose domain is the family of left open intervals in 𝑅 𝑚 is
called an interval function. To distinguish between set function and ordinary functions
defined at each point of a set, we call the latter point functions.

SET OF FIRST SPECIES: A set is said to be a set of first species if it has only a finite
number of derived sets.
SET OF SECOND SPECIES: A set is said to be a set of second species if it has infinite
number of derived sets.
SET OF FIRST SPECIES: A set is said to be a set of first species if it has only a finite
number of derived sets.
SETS: G.Cantor defined a set as a collection of objects of our intuition or thought, within
a certain realm, taken as a whole. Each object in the collection is called an element (or
member) of the set. The notation 𝑎 𝜖 𝐴 means that 𝑎 is an element of the set 𝐴. In this
case we say that 𝑎 is a member of 𝐴 or 𝑎 belongs to 𝐴. The negation of 𝑎 ∇ 𝐴 is written
𝑎 ∈ 𝐴. The set having no element, namely the set 𝐴 such that 𝑎 ∈ 𝐴 for every object 𝑎, is
called the empty set (or null set) and is usually denote by ∅. Two sets 𝐴 and 𝐵 are
identical, i.e., 𝐴 = 𝐵, if every element of 𝐴 belongs to 𝐵, and vice versa. The set
containing 𝑎, 𝑏, 𝑐, … as its elements is said to consist of 𝑎, 𝑏, 𝑐, … and is denoted by
𝑎, 𝑏, 𝑐, … .

A set is called a finite set or an infinite set according as the number of its elements is
finite or infinite.
𝐴 set 𝐴 is a subset of a set 𝐵 if each element of 𝐴 is an element of 𝐵. in this case we also
we also say that 𝐴 is contained in 𝐵 or that 𝐵 contains 𝐴, and we write 𝐴 ⊂ 𝐵. The
negation of 𝐴 ⊂ 𝐵 is 𝐴 ⊄ 𝐵. For every set 𝐴, ∅ ⊂ 𝐴. 𝐴 ⊂ 𝐵 and 𝐵 ⊂ 𝐶 imply 𝐴 ⊂ 𝐶. If
𝐴 ⊂ 𝐵 and 𝐵 ⊂ 𝐴, then 𝐴 = 𝐵. 𝐴 is a proper subset of 𝐵 (in symbols: 𝐴 ⊂
≠
𝐵, if 𝐴 ⊂ 𝐵 and
𝐴 ≠ 𝐵.

SHAPIRO INEQUALITY: Suppose 𝑛 is a natural number and are positive

numbers and:

 𝑛 is even and less than or equal to 12, or

 𝑛 is odd and less than or equal to 23.

Then the Shapiro inequality states that

where .

SHORTEST DISTANCE BETWEEN TWO NON-INTERSECTING LINES: Let 𝒓 = 𝒂 + 𝒕𝒃

………. 1 and 𝒓 = 𝒂’ + 𝒔𝒃’
………. 2

be two non- intersecting lines passing through the points 𝒂, 𝒂’ and parallel to the vector
𝒃, 𝒃’ respectively.

Let 𝑷𝑷′be the shortest distance between the two lines do that 𝑷𝑷′’ is perpendicular to
both b and b’. Then

′
𝒂 − 𝒂′ ∙ 𝒏 𝒂 − 𝒂′ ∙ 𝒃 × 𝒃′ 𝒂 − 𝒂′ , 𝒃, 𝒃′ 𝒂 − 𝒂′ , 𝒃, 𝒃′
𝑷𝑷 = = = = .
𝒏 𝒏 𝒏 𝒃 × 𝒃′

The shortest distance between two non-intersecting lines parallel respectively to 𝒃 and
𝒃’ will be the projection of the line joining any two points – one on each line- on the
vector 𝒃 × 𝒃′ which is perpendicular to both the lines.

If the straight lines intersect, then the shortest distance is zero, the condition for which
is 𝒂 − 𝒂′ , 𝒃, 𝒃′ =0
SHEARING STRAIN: It is defined in terms of the change in angle between two linear
elements from the unstrained shape to the strained state.

SIGNATURE AND INDEX OF A REAL QUADRATIC FORM: Let 𝑦12 + ⋯ + 𝑦𝑝2 − 𝑦𝑝+1
2
− ⋯−
𝑦𝑟2 be a normal form of a real quadratic form X ′ AX of rank r . The number p of positive
terms in a normal form of X ′ AX is called the index of the quadratic form. The excess of
the number of positive terms over the number of negative terms in a normal form of
X ′ AX i.e., p − r − p = 2p − r is called the signature of the quadratic form and is usually
denoted by s.

Thus s = 2p − r.

SIGNATURE OF A PERMUTATION: Let 𝑋 be a finite set, and let 𝐺 be the group of

permutations of 𝑋. There exists a unique homomorphism 𝜒 from 𝐺 to the multiplicative
group {−1, 1} such that 𝜒(𝑡) = −1 for any transposition 𝑡 ∇ 𝐺. The value 𝜒(𝑔), for any
𝑔 ∇ 𝐺, is called the signature or sign of the permutation 𝑔. If 𝜒(𝑔) = 1, 𝑔 is said to be
of even parity; if 𝜒 𝑔 = −1, 𝑔 is said to be of odd parity.
SIGNED MEASURE: Let 𝑋 be a set, and 𝑅 be a ς-ring. A real (complex) valued function 𝜈
on 𝑅 is called a signed measure if it is countably additive as follows: for any 𝐴𝑘 ∇ 𝑅 the
identity 𝐴 = 𝑘 𝐴𝑘 implies the series 𝑘 𝜈(𝐴𝑘 ) is absolute convergent and has the sum
𝜈(𝐴). Any real signed measure 𝜈 has a representation 𝜈 = µ1 − µ2 (𝜈 = µ1 − µ2 + 𝑖µ3 −
𝑖µ4 ), where µ𝑘 are 𝜍 −additive measures.

SIMILARITY OF MATRICES: Let 𝐴 and 𝐵 be square matrices of order 𝑛. Then 𝐵 is said to

be similar to 𝐴 if there exists a non-singular matrix 𝑃 such that 𝐵 = 𝑃−1 𝐴𝑃.

 Similarity of matrices is an equivalence relation.

 Similar matrices have the same determinant.
 Similar matrices have the same characteristic polynomial and hence the same
eigenvalues. If 𝑋 is an eigenvector of 𝐴 corresponding to the eigenvalues λ, then
𝑃−1 𝑋 is an eigenvector of 𝐵 corresponding to the eigenvalue λ where
𝐵 = 𝑃−1 𝐴𝑃.
 If 𝐴 is similar to a diagonal matrix 𝐷, the diagonal elements of 𝐷 are the
eigenvalues of 𝐴.
SIMPLE CIRCUIT: A circuit 𝑣0 𝑣1 𝑣2 . . . 𝑣𝑛 −1 𝑣0 in a graph is said to be simple if the
vertices 𝑣0 , 𝑣1 , 𝑣2 , . . . , 𝑣𝑛−1 are distinct.

SIMPLE FUNCTION: A measurable function 𝑓: 𝑋 → ℝ is simple if it attains only a

∞
countable number of values. A function 𝑓: 𝑋 → ℝ is simple if and only if 𝑓 = 𝑘=1 𝑡𝑘 𝜒𝐴𝑘

for some 𝑡𝑘 ⊆ ℝ and 𝐴𝑘 ∇ 𝐿. That is, simple functions are linear combinations of
indicator functions of measurable sets.

SIMPLE GROUPS: A non-trivial group 𝐺 is said to be simple if the only normal subgroups
of 𝐺 are the whole of 𝐺 and the trivial subgroup {𝑒} whose only element is the identity
element e of 𝐺.
SIMPLE PATH: A simple path in a graph is a path that contains no vertex more than once.
By definition, cycles are particular instances of simple paths.
SIMPLE-PERIODIC FUNCTION: A periodic function which has only one fundamental
period is said to be simply-periodic.

SIMPLEX ALGORITHM: The various steps involved in the computation of optimal

solution are as follows:

Step 1: First of all, it should be checked whether the given problem is to be maximized
or minimized. If the given problem is of minimization, convert it into minimization by
multiplying both sides of the objective function by (-1).

i.e. Minimization z = - Maximum (-z)

Step 2: Make all 𝑏𝑖 ’s i = 1 ,2,……….m non –negative , i.e., the RHS of each of the
constraints should be non- negative. This can done by multiplying with (-1) to that
inequation which contains negative 𝑏𝑖 .

Step 3: Convert inequalities of constraints of constraints into equations. For this step
introduce slack and surplus variables and the coefficients of these variables must be
equated to zero in the objective function.

All these variables must be introduced if the objective function is of maximization. If it is

not so, then the objective function must be converted to maximization. It must be kept
in mind that all the imaginary variables added to the constraints will form a part of the
category basic variables if they form an identity matrix.
Step 4: To find initial feasible solution. Write the constants in matrix form. The variables
corresponding to the basic columns are the basic variables, other are non basic
variables , setting each non- basic variables equal to zero, and using constants equations
, we get the initial basic feasible solution.

Step 5: Construction of starting simplex table.

𝑐𝑗 𝑐1 𝑐2 𝑐3 𝑐4 etc.
𝑋𝑗 𝑋1 𝑋2 𝑋3 𝑋4 etc .
Basic 𝐶𝐵 𝑋𝐵 𝐴1 𝐴2 𝐴3 𝐴4 etc . Minimum
variables ratio

The table can be constructed in a systematic ways as follows:

I. At the top of the table, write the numerical values of the constants (𝑐𝑖 ) and x-
variables 𝑥1 to be obtained from the objective function.
II. 𝐴1 𝐴2 , … denotes the columns of the coefficients matrix A are to be obtained from
the matrix equations of the constraints. Write the columns of coefficients of
𝑥1 below 𝐴1 column of coefficient of 𝑥2 below 𝐴2 and so on.
III. The variable 𝑥𝑗 above the basic unit column are the basic variables. These basic
variables to be written in the column of basic variables (the left column of the
table)
IV. 𝐶𝑏 is the column of numerical of the cost the basic variables. Write the cost
entries below 𝐶𝐵 according to the x-variables in the left column.
V. 𝑋𝐵 is the column at numerical vales of the best variables. Variables Entries of this
column at initial stage are nothing but the value 𝑏1 , 𝑏2 …given in the constraints.
At this stage is called initial basic feasible solution.

To test whether the solution given in the tables (IBFS) is optimal we first calculate the
net evaluation ∆𝑗 = 𝑧𝑗 − 𝑐𝑗 = 𝐶𝑗 𝐴𝑗 − 𝑐𝑗 each j. Then we apply the following optimal test.

Step 6: Calculation of Net Evalutaion. ∆𝑗 = 𝑧𝑗 − 𝑐𝑗 = 𝐶𝑗 𝐴𝑗 − 𝑐𝑗

For 𝑗 = 1, ∆1 = 𝐶𝐵 𝐴1 − 𝑐𝑗

Here , 𝐶𝐵 is the column vector give in the table at second column 𝐴1 is the first column of
the coefficient matrix, which as also given in the table 𝐶1 is the cost of 𝑥1 given at the top
of the table.

Then we write 𝐶𝐵 and 𝐴1 as pairs of tuples. The product 𝑧1 = 𝐶𝐵 𝐴1 is given by the sum of
product of their corresponding coordinates. Subtracting the value 𝐶1 from it to get
∆1 = 𝑍1 − 𝑐𝑗 .

In a similar way we can find ∆1 , ∆3 , …, and so on

Write all ∆𝑗 in the bottom of the table.

Step 7: Optimality Test. Here, we apply the following optimality test:

(i) If all ∆𝑗 ≥ 0, then solution is optimal.

(ii) IF all ∆𝑗 ≥ 0, 𝑎𝑛𝑑 ∆𝑗 > 0 corresponding to at least one non-basic variable, then
optimal solution is unique.
(iii) All ∆𝑗 ≥ 0, 𝑎𝑛𝑑 ∆𝑗 = 0 corresponding to at least one non-basic variable, then
optimal solution is not unique.
(iv) If ∆𝑗 < 0 for the least one 𝑗, then solution is not optimal, and required further
improvement.

If solution is not optimal, then we apply the following steps:

Step 8: Incoming Column Vector and Incoming Variable. Select the most negative
∆𝑗 = 𝑧𝑗 − 𝑐𝑗 say ∆, is most negative. Express it by putting it in a circle. The corresponding
column 𝐴𝑟 is the incoming column (it is to be converted to some basic unit column) and
the corresponding variable 𝑥𝑟 (at the top above the O) is incoming variable (i.e. it will
become a basic variable and will replace some basic variable presently appeared in the
left column of basic variables)

Step 9: Minimum Ratio. We calculate the ratios of elements of 𝑋𝐵 and incoming column
vector 𝐴𝑟 and written the last column of the table. The ratio which are negative or
undefined (infinity) are to be ignored. Out of the remaining ratios, the minimum ratio is
selected. Express it by putting a circle O on this minimum ratio.
Step 10: Pivoting Element. It is known as key element ratio also. The element of the
incoming column vector 𝐴𝑟 appears at the row of the minimum ratio is called pivot or
key element. Express it in the table by putting . Convert the key element to 1 and
other entries of this column shall be converted to zero, so that the declared incoming
column vector 𝐴𝑟 will become one of basic (unit) column vector.

Step 11: Outgoing Column Vector and Basic Variable. Select the basic column whose unit
element 1 appears in this row of key element. That basic column vector is to be declared
outgoing column vector. Express it by a downwards arrow below it. The variable
corresponding to the outgoing vector is the outgoing basic variable.

Step 12: Improved Simplex Table: Using the starting table, pick up the matrix
𝑋𝐵 , 𝐴 = 𝑋𝐵 , 𝐴1 , 𝐴2 , …

By applying row operations convert the key element to 1 and other elements of its
column 𝐴𝑟 to 0. Then A, becomes a basic column.

From this row reduced matrix, take the change values back to the table for the columns
𝑋𝐵 , 𝐴1 , 𝐴2 , …. and so on and construct an improved simplex table.

Repeat the steps (6) and (7) on this table. Here, if solution is not optimal, continue the
steps for improved solution.

SIMPLEX METHOD: Simplex Method is one of the most efficient and convenient method
of solution of problems involving more than two variables. The simplex algorithm
proceeds in a step- by step systematic manner and it involves feasible solutions which
are provided by the corner points only. This method is also used to indicate whether a
given solution is optimal or not.

Initialize Set up slack variables, which are the initial basic variables. Set
decision variables to be the initial non-basic variables.
Optimality Test Basic feasible solution is optimal iff every coefficient in row 0 is
non-negative. If so, we are done, otherwise we will continue.
Iteration Settle on entering basic variable – having the most negative
coefficient in row 0. This is the pivot column. Settle on leaving basic
 variable by applying the minimum ratio test.
SIMPLY AND MULTIPLY CONNECTED DOMAINS: A domain in which every closed curve
can be shrunk to a point without passing out of the region, is called a simply connected
domain. If a domain is not simply connected, then it is called multiply connected
domain.

SIMULATION: This is a numerical experiment in a simulated model of a phenomenon

which we want to analyze. Simulation is one of the most popular techniques in OR.
Singleton: 1) In object-oriented programming , a singleton class is a class that can have
only one object(an instance of the class) at a time. For example, using Visual C++ , the
"Application" class is an example of a singleton class. You can only create only one
object of an Application class.

2) In set theory (mathematics), a singleton is a set with a single element. An example is

the set S of all integer s that are neither positive nor negative. In this case: S = {0}

3) In certain card games, a singleton is a card that is the only one of a suit in a hand. In
the game of bridge, each player is dealt 13 cards and any hand with only one card in a
given suit is said to hold a singleton of that suit (which, depending on the play of the
cards may turn out to be a strength or a weakness, or make no difference.)

SINGULAR ELEMENTS: Elements which are not regular in a Banach algebra A are called
singular elements. The element 0 (null vector) is always singular in A. The set of
singular elements in A is denoted by S.
SINGULAR POINT (COMPLEX ANALYSIS): A point 𝑧 = 𝑧0 is said to be a singular point of
a function 𝑓 𝑧 if 𝑓′(𝑧0 ) does not exist.

SINGULAR VALUE DECOMPOSITION: If T is any compact operator on a separable Hilbert

space then there exists orthonormal sequences (𝑒𝑛 ) and (𝑓𝑛 ) such that 𝑇𝑥 =
𝑘 µ𝑘 ⟨ 𝑥, (𝑒𝑛 ) ⟩ (𝑓𝑛 ) where (µ𝑘 ) is a sequence of positive numbers such that µ𝑘 → 0 if it
is an infinite sequence.

Every finite-dimensional normed vector space is a Banach space.

SINK: A source of negative strength −𝑚 or inward radial flow is called a sink. This is a
point at which the fluid continuously annihilates. Whirlpool is an example of a sink.
SKEW- HERMITIAN MATRIX: A square matrix A = aij is said to be skew-Hermitian if
th th
the i, j element of A is equal to the negative of the conjugate complex of the j, i
element of A i.e., if aij = aji for all i and j.

If A is a skew-Hermitian matrix, then aii = aii by definition . ∴ aii = aji =0 i.e., aii must
be either a pure imaginary number or must be zero.Thus the diagonal elements of a
skew-Hermitian matrix must be pure imaginary number or zero.

0 −2 − i −i 3 + 4i
Illustration: , are skew-Hermitian matrices. A skew-
2−i 0 −3 + 4i 0
Hermitian matrix over the field of real numbers is nothing but a real skew-symmetric
matrix. Obviously, a necessary and sufficient condition for a matrix A to be skew-
Hermitian is that A° = −A.

SKEWNESS: Skewness is asymmetry in a statistical distribution, in which the curve

appears distorted or skewed either to the left or to the right. Skewness can be
quantified to define the extent to which a distribution differs from a normal
distribution. In a normal distribution, the graph appears as a classical, symmetrical
"bell-shaped curve." The mean, or average, and the mode, or maximum point on the
curve, are equal.

 In a perfect normal distribution (green solid curve in the illustration below), the
tails on either side of the curve are exact mirror images of each other.
 When a distribution is skewed to the left (red dashed curve), the tail on the curve's
left-hand side is longer than the tail on the right-hand side, and the mean is less than
the mode. This situation is also called negative skewness.
 When a distribution is skewed to the right (blue dotted curve), the tail on the
curve's right-hand side is longer than the tail on the left-hand side, and the mean is
greater than the mode. This situation is also called positive skewness.
SKEW- SYMMETRIC MATRIX: A square matrix A = aij is said to be skew-symmetric if
th
the i, j element A i.e., if aij = aij for all i, j.

If A is a skew-symmetric matrix, then aij = −aji be definition ∴ aii = −aii , for all
values of i ∴ 2aii = 0 or aii = 0

Thus the diagonal elements of a skew-symmetric matrix are all zero.

 0 h g 0 −3i −4
Illustration: The matrices −h 0 f , 3i 0 8 are skew-symmetric matrices.
−g −f 0 4 −8 0

A necessary and sufficient condition for a matrix A to be skew-symmetric is that

𝐴′ = – 𝐴.

SLACK VARIABLES: If the given constrains in LPP has a sign ≤ , then to turn the
inequations as equality a positive quantity is to be added on LHS. The positive quantities
(or variables) which are added on the LHS are called slack variables.

SLATER’S WEAK CONSTRAINT QUALIFICATION: Let X 0 be an open set in Rn and let 𝑔 be

an 𝑚- dimensional vector function defined on X 0 , and let 𝑋 = x: x ∇ X 0 , g(x) ≤ 0 . 𝑔 is
said to satisfy the Slater’s weak constraint qualification at x ∇ X if 𝑔 is differentiable at
x, 𝑔𝐼 is pseudoconvex at x, and there exists an x ∇ X such that

g I x > 0 where

𝐼 = i: g i x = 0

SMOOTH MAP: Let 𝑋 ⊂ 𝑅 𝑛 and 𝑌 ⊂ 𝑅 𝑙 be arbitrary subsets. A map 𝑓: 𝑋 → 𝑌 is said to

be smooth at 𝑝 ∇ 𝑋, if there exists an open set 𝑊 ⊂ 𝑅 𝑛 around p and a smooth map
𝐹: 𝑊 → 𝑅 𝑙 which coincides with 𝑓 on 𝑊 ∩ 𝑋. The map 𝑓 is called smooth if it is
smooth at every 𝑝 ∇ 𝑋. If f is a bijection of 𝑋 onto 𝑌 , and if both 𝑓 and 𝑓 −1 are smooth,
then f is called a diffeomorphism. A smooth map 𝐹 as above is called a local smooth
extension of 𝑓. In order to show that a map defined on a subset of 𝑅 𝑛 is smooth, one
thus has to find such a local smooth extension near every point in the domain of
definition. It is easily seen that a smooth function is continuous
SMOOTH VECTOR FIELD: A smooth vector field on a manifold 𝑆 ⊂ 𝑅 𝑛 is a smooth map
𝑌 ∶ 𝑆 → 𝑅 𝑛 such that 𝑌 (𝑝) ∇ 𝑇𝑝 𝑆 for all 𝑝 ∇ 𝑆.
S-NORM IN FUZZY SETS: Any function 𝑠 ∶ [0, 1] × [0, 1] → [0, 1] that satisfies
following Axioms 𝑠𝑙 − 𝑠4 is called an 𝑠 −norm.
Axiom sl: 𝑠(1, 𝑙) = 1, 𝑠(0, 𝑎) = 𝑠(𝑎, 0) = 𝑎 (boundary condition).
Axiom s2: 𝑠(𝑎, 𝑏) = 𝑠(𝑏, 𝑎) (commutative condition).
Axiom s3: If 𝑎 ≤ 𝑎′ and 𝑏 ≤ 𝑏′, then 𝑠(𝑎, 𝑏) ≤ 𝑠(𝑎′ , 𝑏′) (nondecreasing condition).
Axiom s4: 𝑠(𝑠(𝑎, 𝑏), 𝑐) = 𝑠(𝑎, 𝑠(𝑏, 𝑐)) (associative condition).
SOLUTION OF AN LPP: The solution of a LPP is that set of values of the assigned
variables which satisfies the given constraints of a LPP.

SOLUION OF TRAVELLING SALESMAN PROBLEM: The problem can be solved by

assignment technique. The solution thus obtained may not be feasible for this problem.
For example, if we choose the elements 𝑐15 , 𝑐23 , 𝑐34, 𝑐42 , 𝑐51 , then he corresponding
solution gives 𝐴1 → 𝐴5 , 𝐴5 → 𝐴1 , i.e., 𝐴1 is followed by 𝐴5 and𝐴5 𝑏𝑦 𝐴1 . This violates
our restriction. In such cases after solving the given problem by assignment technique,
we use the method of enumeration by assigning the next minimum element of the
matrix in place of zero.

Such problems occur in the field of postal deliveries, school bus routing, television
relays assembly lines, production of several items by one machine.

SOLVABLE GROUPS: The concept of a solvable group was introduced into mathematics
by Evariste Galois, in order to state and prove his fundamental general theorems
concerning the solvability of polynomial equations. A group 𝐺 is said to be solvable (or
soluble) if there exists a finite sequence 𝐺0 , 𝐺1 , . . . , 𝐺𝑛 of subgroups of 𝐺, where 𝐺0 = {1}
and 𝐺𝑛 = 𝐺, such that 𝐺𝑖−1 is normal in 𝐺𝑖 and 𝐺𝑖 /𝐺𝑖−1 is Abelian for 𝑖 = 1, 2, . . . , 𝑛.

SOLVABLE POLYNOMIALS: We say that a polynomial with coefficients in a given field is

solvable by radicals if the roots of the polynomial in a splitting field can be constructed
from its coefficients in a finite number of steps involving only the operations of addition,
subtraction, multiplication, division and extraction of nth roots for appropriate natural
numbers 𝑛.

SOURCE: Any point in a two dimensional field is assumed to flow uniformly in all
directions along radial paths or consists of outward radial flow from a point, is called a
simple source. The lines of flow will be straight radial lines i.e., q θ = 0 and q r (r). A
source is a point at which the fluid is continuously created. In fact, this is a purely
abstract conception which foes not occur in nature. If the total flux outwards across a
small closed surface surrounding the point be +𝑚 then +𝑚 is called the strength of
the source.

SPACE OF MEASURABLE FUNCTIONS: Let (𝑋, 𝐿, µ) be a measure space. For 1 ≤ 𝑝 < ∞,

we define 𝐿𝑝 (µ) to be the space of measurable functions 𝑓: 𝑋 → 𝐾 such that
 1
∫𝑋 𝑓 𝑝 𝑑𝜇 < ∞. We define || · ||𝑝 ∶ 𝐿𝑝 (µ) → [0, ∞) by 𝑓 𝑝 = [∫𝑋 𝑓 𝑝 𝑑𝜇] 𝑝. Notice that
if 𝑓 = 0 almost everywhere, then | 𝑓 |𝑝 = 0 almost everywhere, and so ||𝑓||𝑝 = 0.
However, there can be non-zero functions such that 𝑓 = 0 almost everywhere. So || · ||𝑝
is not a norm on 𝐿𝑝 (µ).

SPANNING TREE: A spanning tree is a connected network of the 𝑛 nodes that contains
no undirected cycles. A network with 𝑛 nodes requires only ( 𝑛 – 1 ) links to provide a
path between each pair of nodes. Every connected graph contains a spanning tree.
SPACE: Space is a term that can refer to various phenomena in science, mathematics,
and communications. In astronomy and cosmology, space is the vast 3-dimensional
region that begins where the earth's atmosphere ends. Space is usually thought to begin
at the lowest altitude at which satellites can maintain orbits for a reasonable time
without falling into the atmosphere. This is approximately 160 kilometers (100 miles)
above the surface. Astronomers may speak of interplanetary space (the space between
planets in our solar system), interstellar space (the space between stars in our galaxy),
or intergalactic space (the space between galaxies in the universe). Some scientists
believe that space extends infinitely far in all directions, while others believe that space
is finite but unbounded, just as the 2-space surface of the earth has finite area yet no
beginning nor end.

In mathematics, space is an unbounded continuum (unbroken set of points) in which

exactly three numerical coordinates are necessary to uniquely define the location of any
particular point. It is sometimes called 3-space because it contains three distance
dimensions. If a continuum requires fewer or more than three coordinates
(dimensions) to uniquely define the location of a point, that continuum is sometimes
called n-space or n-dimensional space, where n is the number of dimensions. Thus, for
example, a line constitutes 1-space and a plane constitutes 2-space. When time is
considered as a dimension along with the usual three in conventional space, the result is
sometimes called 4-space, 4-dimensional space, time-space, or space-time. In digital
communications, the term space refers to an interval during which no signal is
transmitted, or during which the signal represents logic 0. The term space may also be
used in reference to the time interval separating two characters, bytes, octets, or words
in a digital signal.
SPACE CURVE: A Curve in Eucliden space of three dimension is the locus of a point
whose position vector 𝑟 with respect to origin say 𝑂 is function of single parameter 𝑡.
The Cartesian coordinates (𝑋, 𝑌, 𝑍) of point are called components of 𝑅 and are function
of parameter. Therefore we can express the equation of curve in Thus 𝑟 𝑡 = 𝑥 𝑡 𝑖 +
𝑦 𝑡 𝑗 + 𝑧 𝑡 𝑘 or simply 𝑅 = (𝑥, 𝑦, 𝑧) represents a curve in the space.

The curve is known as plane curve if it lies on a plane, otherwise it is said to be a skew
twisted or tortuous curve.
The parametric equations of a curve are
𝑥 = 𝑥 𝑡 , 𝑦 = 𝑦 𝑡 , 𝑧 = 𝑧( 𝑡)
Where x, y, z are real valued functions of a single real parameter 𝑡 ranging over a set of
vaues 𝑎 ≤ 𝑡 ≤ 𝑏.
SPANNING VECTORS: The vectors 𝑣1 , 𝑣2 , ⋯ ⋯ , 𝑣𝑛 in 𝑉 span 𝑉 if every vector v 𝜖 𝑉 is a
linear combination 𝛼1 𝑣1 + 𝛼2 𝑣2 + ⋯ ⋯ + 𝛼𝑛 𝑣𝑛 of 𝑣1 , 𝑣2 , ⋯ ⋯ , 𝑣𝑛 .

SPEARMAN’S RANK CORRELATION: Set 𝑑𝑖 = 𝑅𝑖 − 𝑆𝑖 . Then

𝑑𝑖2
𝑟𝑠 = 1 − 6
𝑛3 − 𝑛

is called Spearman’s rank correlation. If there is no dependence between 𝑋 and 𝑌, i.e., if

the 𝑋𝑖 and 𝑌𝑗 are independent random variables, then 𝐸 𝑟𝑠 = 0 𝑎𝑛𝑑 𝑉 𝑟𝑠 = (𝑛 − 1)−1 .

SPECTRAL RADIUS: The spectral radius of T is 𝑟(𝑇) = 𝑠𝑢𝑝{ λ : λ ∇ ς(T)}.

Note that
1
1. For a bounded operator T, we have 𝑟(𝑇) = lim𝑛 →∞ 𝑇 𝑛 𝑛.

2. There exists 𝜆 ∇ 𝜍(𝑇) such that | 𝜆 | = 𝑟(𝑇).

3. For any 𝑇 ∇ 𝐵(𝐻) the spectrum of the adjoint operator is 𝜍(𝑇 ∗ ) = {𝜆: 𝜆 ∇ 𝜍(𝑇)}.
4. If 𝑈 is a unitary operator then 𝜍(𝑈) ⊆ {| 𝑧 | = 1}.
5. If 𝑇 is Hermitian then 𝜍(𝑇) ⊆ ℝ.
6. The operator 𝑈: 𝑓(𝑡) ↦ 𝑒 𝑖𝑡 𝑓(𝑡) on 𝐿2 [0,2𝜋] is unitary and has the entire unit
circle {| 𝑧 | = 1} as its spectrum .

SPECTRAL THEOREM: Let 𝑇 be an arbitrary operator on 𝐻, and 𝑙𝑒𝑡 𝜆1 𝜆2 , … , 𝜆𝑛 be the

distinct eigen values of 𝑇 with corresponding eigen spaces 𝑀1 , 𝑀2 , … , 𝑀𝑛 . Let 𝑃1 𝑃2 , … , 𝑃𝑛
 be the perpendicular projections respectively on these eigen spaces. Then the following
statements are equivalent to one another:

i. The subspaces 𝑀𝑖 ′ 𝑠 are pariwise orthogonal and span 𝐻.

𝑛
ii. 𝑃𝑖 ′ 𝑠 are pariwise orthogonal 𝑖=1 𝑃𝑖 = 𝐼, 𝑎𝑛𝑑
𝑛
𝜆𝑖 𝑃𝑖 = 𝑇.
𝑖=1

iii. 𝑇 is a normal operator.

SPECTRAL THEOREM: For every self-adjoint operator 𝐻in a Hilbert space 𝑋, there
exists a unique real spectral measure 𝐸 such that

∞
𝐻 = ∫−∞ 𝜆𝐸(𝑑𝜆).

In other words, 𝐻 and 𝐸 correspond to each other by the relations

∞
𝐷 𝐻 = 𝑥 ∫−∞ 𝜆2 𝐸 𝑑𝜆 𝑥, 𝑥 < +∞ ,

∞
𝐻𝑥, 𝑦 = ∫−∞ 𝜆(𝐸 𝑑𝜆 𝑥, 𝑦),

𝑥 ∇ 𝐷, 𝑦 ∇ 𝑋.

This is the spectral theorem for self-adjoint operators. The support of 𝐸 is equal to the
spectrum 𝜍(𝐻), so that we can write

∞
𝐻 = ∫𝜍(𝐻) 𝜆𝐸(𝑑𝜆) = ∫−∞ 𝜆𝜒𝜍 𝐻 (𝜆)𝐸(𝑑𝜆),

Where 𝜒𝑚 stands for the characteristic function of 𝑀.

SPECTRAL THEOREM FOR COMPACT NORMAL OPERATORS: Let 𝑇 be a compact normal

operator on a Hilbert space H. Then there exists an orthonormal sequence (𝑒𝑛 ) of
eigenvectors of T and corresponding eigenvalues (𝜆𝑛 ) such that:
𝑇 𝑥 = 𝑛 𝜆𝑛 x, 𝑒𝑛 𝑒𝑛 for all x ∇ H. (1)

If λn) is an infinite sequence it tends to zero.

Conversely, if T is given by a formula (1) then it is compact and normal.

Let T be a compact normal operator then

 1. The set of of eigen values of 𝑇 is either finite or a countable sequence tending to
zero.

2. All the eigenspaces, i.e. 𝑘𝑒𝑟(𝑇 − 𝜆 𝐼), are finite-dimensional for all 𝜆 ≠ 0.

This Lemma is true for any compact operator.

Note that

1. Let T be a compact normal operator. Then all non-zero points 𝜆 ∇ 𝜍(𝑇) are
eigenvalues and there exists an eigenvalue of modulus ||𝑇||.
2. Let T be a compact normal operator on a separable Hilbert space 𝐻, then there
exists a orthonormal basis 𝑔𝑘 such that 𝑇 𝑥 = 𝑛 𝜆𝑛 x, 𝑔𝑛 𝑔𝑛 and 𝜆𝑛 are
eigenvalues of T including zeros.

SPECTRAL THEOREM FOR COMPACT SELF-ADJOINT OPERATORS IN HILBERT SPACE:

Let 𝑇 be a compact self-adjoint operator in a Hilbert space 𝑋. Then there is an
orthonormal basis consisting of eigenvectors of 𝑇.

SPHERE: Sphere can be regarded as a surface of revolution formed by rotating the circle
𝑥² + 𝑧² = 𝑎², 𝑦 = 0 [in zx-plane] about z-axis.

Co-ordinates of any point, say Q on any point of the circle can be taken as 𝑥 =
𝑎 sin 𝑢, 𝑦 = 0, 𝑧 = 𝑎 cos 𝑢 whence the equation of the sphere can be written as

𝑥 = 𝑎 sin 𝑢 𝑐𝑜𝑠𝜙, 𝑦 = 𝑎 sin 𝑢 sin 𝜙, 𝑧 = 𝑎 cos 𝑢

or 𝑥 = 𝑎 sin 𝑢 cos 𝑣, 𝑦 = 𝑎 sin 𝑢 sin 𝑣, 𝑧 = 𝑎 cos 𝑢;

Where 𝜙 is replaced by v,

or in vector form 𝑟 = (𝑎 sin 𝑢 cos 𝑣, 𝑎 sin 𝑢 sin 𝑣, 𝑎 𝑐𝑜𝑠𝑢).

SPHERICAL HARMONICS: In dealing with the theory of potential, we commonly use the
Laplace’s equation

∂2 V ∂2 V ∂2 V
+ + =0
∂x 2 ∂y 2 ∂z 2

Any solution Vn of this equation, which is a homogeneous polynomial (of degree n) in

x, y, z is called the Solid Spherical Hermonics.
Laplace’s equation in polar or spherical coordinates is written as
∂2 V 2 ∂V 1 ∂2 V 𝑐𝑜𝑡𝜃 𝜕𝑉 1 𝜕2𝑉
+ + + + =0
∂r 2 r ∂r r 2 ∂θ2 𝑟 2 𝜕𝜃 𝑟 2 𝑠𝑖𝑛2 𝜃 𝜕∅2

Let Vn ; when expressed in polar coordinates take the form Vn = r n . Un (θ, ∅), where
Un (𝜃, ∅) is a function of θ and ∅ only.

Then Un (θ, ∅) is called the surface Spherical Harmonics of degree n.

Obviously when r = l, Vn = Un .

Hence surface spherical harmonics is the value of the corresponding solids spherical
harmonics on the surface of a unit sphere whose centre is the origin.

V
Kelvin’s Theorem:- If Vn is a solid harmonics of degree n, then r 2nn+1 is a solid spherical

harmonics of degree −(n + 1).

SPHERICAL INDICATRICES: When we move all unit tangent vectors 𝑡 of a curve 𝐶 to a

point, their extremities describes a curve 𝐶1 on the unit sphere; this curve 𝐶1 is called the
spherical image (spherical indicatrix) of 𝐶. There is a one-to-one correspondence
between 𝐶 and 𝐶1 . We can similarly obtain the spherical image of 𝐶 when its bonormals
or principal normals are moved to a point.
SPHERICAL NEIGHBOURHOOD OF A: A spherical neighbourhood of the point 𝑓 of radius
𝑟 is the set of points.

𝑆𝑟 𝑓 = 𝑔 𝜖 𝑋 ∶ 𝜌 𝑓, 𝑔 < 𝑟 .

If the metric space is 𝑅 𝑛 , 𝜌) then

𝑆𝑟 𝑓 = 𝑔 𝜖 𝑅 𝑛 ∶ 𝑔 − 𝑓 < 𝑟

SPHERICAL TRIANGLES: A spherical triangle ∆ consists of three vertices 𝐴, 𝐵, 𝐶 ∇ 𝑆 2

and three geodesics from 𝐵 to 𝐶, from 𝐶 to 𝐴, and from 𝐴 to 𝐵. The geodesics are the
sides of ∆ and their lengths will be denoted by 𝑎, 𝑏, and 𝑐 respectively. If 𝛾 is a geodesics
from 𝐴 to 𝐵, then 𝛾 is of the form 𝛾 ∶ [0, 𝑐] → 𝑆 2 ; 𝑡 → (𝑐𝑜𝑠𝑡)𝐴 + (𝑠𝑖𝑛 𝑡)𝑢 where 𝑢 is
the unit tangent vector to the side 𝐴𝐵 of ∆ at 𝐴. Similarly, the unit tangent vector to the
side 𝐴𝐶 at 𝐴 is a vector 𝑣. The angle between these unit vectors 𝑢 and 𝑣 is called the
angle between the two sides 𝐴𝐵 and 𝐴𝐶 and is denoted by 𝛼. Similarly, we denote the
angle between the two sides 𝐵𝐴 and 𝐵𝐶 at 𝐵 is denoted by 𝛽 and the angle between the
two sides 𝐶𝐴 and 𝐶𝐵 at 𝐶 is denoted by 𝛾.

SPHEROID: A spheroid is similar to a sphere but is lengthened or shortened in one

dimension.
SPIRAL: The curve 𝑟 = 𝑎𝜃, graphed in polar coordinates, has a spiral shape.

SPLITTING FIELD: Let 𝐿: 𝐾 be a field extension, and let 𝑓 ∇ 𝐾[𝑥] be a polynomial with
coefficients in 𝐾. The polynomial 𝑓 is said to split over 𝐿 if 𝑓 is a constant polynomial or
if there exist elements 𝛼1 , 𝛼2 , . . . , 𝛼𝑛 of 𝐿 such that 𝑓(𝑥) = 𝑐(𝑥 − 𝛼1 )(𝑥 − 𝛼2 ) · · · (𝑥 −
𝛼𝑛 ), where 𝑐 ∇ 𝐾 is the leading coefficient of 𝑓.
SPORADIC GROUPS: The sporadic groups are the 26 finite simple groups that do not fit
into any of the four infinite families of finite simple groups (i.e., the cyclic groups of
prime order, alternating groups of degree at least five, Lie-type Chevalley groups,
and Lie-type groups). The smallest sporadic group is the Mathieu group , which has
order 7920, and the largest is the monster group, which has
order .
SQUARE FREE INTEGER: An integer 𝑑 is said to be square-free if there is no integer 𝑘
other than ±1 for which 𝑘 2 divides 𝑑.
SQUARE MATRIX: An m × n matrix for which m = n(i.e., the number of rows is equal to
the number of columns) is called a square matrix of order n. It is also called an n-rowed
square matrix. Thus is a square matrix, we have the same number of rows and columns.
The elements 𝑎𝑖𝑗 of a square matrix 𝐴 = 𝑎𝑖𝑗 , for which 𝑖 = 𝑗 i.e., the elements
𝑚 ×𝑛

𝑎11 , 𝑎22 , 𝑎33 , … . , 𝑎𝑚𝑛 are called the diagonal elements and the line along which they lie
is celled the principal diagonal of the matrix.

Example the matrix

0 1 2 3
2 3 1 0
A=
5 0 1 1
0 0 1 2 4×4

is a square matrix of order 4. The elements 0,3,1,2 constitute the principle diagonal of
this matrix.

STABLE SOLUTION: Solution to a problem is stable if a small modification in the

conditions of the problem does not change the solution too much. How much is much
depends of course on the problem.
STANDARD DECODING ARRAY (SDA): A table containing all the pairs (𝑒, 𝑆(𝑒)) where e is
the leader of a coset is called a syndrome lookup table or standard decoding array
(SDA).
STANDING OR STATIONARY WAVES: Let the two simple harmonic progressive wave of
the same amplitude, wave length and period travel in opposite directions are given by
the equations

y1 x, t = A sin mx − nt and y2 x, t = A sin mx − nt

Let y x, t represents the free surface profile then by superimposing these two waves,
we get

y x, t = y1 + y2 A sin mx − nt + sin⁡
(mx + nt)

or y x, t = 2A sin mx cos nt.

A wave of this type is known a standing or stationary wave which is not propagated. At
any instant 𝑡 the form of the surface (2) is sine curve of amplitude 2𝑎 cos 𝑛𝑡 which
varies between 0 to 2𝑎. The points of intersection with 𝑋 - axis are given by

𝜌𝜋
sin 𝑚𝑥 = 0 𝑜𝑟 𝑥 = ; 𝜌 = 0, ±1, ±2
𝑚

Surface waves arise where the vertical acceleration of fluid is no longer neglected and
the wave length is small in comparison with the depth of the water. It moves along the
surface of liquid. The disturbance does not extend far below the surface. Surface waves
occur in deep and bounded ( in horizontal directions) liquids like ocean and lakes.

Tidal waves describe the alternative limit where the wave length of oscillations is much
larger compared to the water. The vertical accelerations can be assumed negligible
compared with the horizontal accelerations. The disturbance affects the motion of the
whole of fluid. Tidal waves are also known as long wave in shallow water.

STATIONARY PROCESSES: 1. A sequence {𝑋𝑡 , 𝑡 ∇ 𝑍} is strongly stationary or strictly

𝐷
stationary if (𝑋𝑡1 , . . . , 𝑋𝑡𝑘 ) (𝑋𝑡1+𝑕 , . . . , 𝑋𝑡𝑘+𝑕 ) for all sets of time points 𝑡1 , . . . , 𝑡𝑘 and
=
integer 𝑕.
2. A sequence is weakly stationary, or second order stationary if
(a) 𝐸(𝑋𝑡 ) = µ, and
(b) cov(𝑋𝑡 , 𝑋𝑡+𝑘 ) = 𝛾𝑘 , where µ is constant and 𝛾𝑘 is independent of 𝑡.
3. The sequence {𝛾𝑘 , 𝑘 ∇ 𝑍} is called the autocovariance function.
STATISTIC: A statistic is a quantity calculated from the items in a sample.
STATISTICAL ANALYSIS: Statistical analysis is a component of data analytics. In the
context of business intelligence (BI), statistical analysis involves collecting and
scrutinizing every data sample in a set of items from which samples can be drawn. A
sample, in statistics, is a representative selection drawn from a total population.
Statistical analysis can be broken down into five discrete steps, as follows:

 Describe the nature of the data to be analyzed.

 Explore the relation of the data to the underlying population.
 Create a model to summarize understanding of how the data relates to the
underlying population.
 Prove (or disprove) the validity of the model.
 Employ predictive analytics to run scenarios that will help guide future actions.
The goal of statistical analysis is to identify trends. A retail business, for example, might
use statistical analysis to find patterns in unstructured and semi-structured customer
data that can be used to create a more positive customer experience and increase sales.

STATISTICAL INFERENCE: Statistical inference refers to the process of estimating

unobservable characteristics on the basis of information that can be observed. The
complete set of all items of interest is called the population. The characteristics of the
population are usually not known. In most cases it is too expensive to survey the entire
population. However, it is possible to obtain information on a group randomly selected
from the population. This group is called a sample. An unknown characteristic of a
population is called a parameter. A quantity that is calculated from a sample is called a
statistic. In many cases the value of a statistic is used as an indicator of the value of a
parameter. This type of statistic is called an estimator.
STATISTICS: Statistics is the study of ways to analyze data. It consists of descriptive
statistics and statistical inference.
STEADY AND UNSTEADY FLOWS: Steady flow occurs when at various points of the flow
field the conditions and properties associated with the fluid flow remain unaltered with
∂A
time. Mathematically, it can be expressed as = 0, where A represents the
∂t

characteristic of the fluid, e.g., velocity, density, temperature and pressure etc. Thus, in
steady motion time drops out of the independent variables and the various field
quantities becomes function of the space coordinates. Water being pumped through a
fixed system at a constant rate is an example of steady flow.

The flow is said to be unsteady when conditions at any point change with regard to the
time, what being pumped through a fixed system at an increasing rate is an example of
unsteady flow.

STEINITZ EXCHANGE LEMMA: If 𝐵 = {𝑏1 , . . . , 𝑏𝑛 } is a basis of a vector space 𝑉 and

𝑛
𝑣 = 𝑗 =1 𝜆𝑗 𝑏𝑗 with 𝜆1 ≠ 0, then {𝑣, 𝑏2 , . . . , 𝑏𝑛 } is a basis of 𝑉 .
STERADIAN: The steradian (symbolized sr) is the Standard International (SI) unit of
solid angular measure. There are 4 pi, or approximately 12.5664, steradians in a
complete sphere. A steradian is defined as conical in shape, as shown in the illustration.
Point P represents the center of the sphere. The solid (conical) angle q, representing
one steradian, is such that the area A of the subtended portion of the sphere is equal
to r2, where r is the radius of the sphere.

A general sense of the steradian can be envisioned by considering a sphere whose

radius is one meter (r = 1m). Imagine a cone with its apex P at the center of the sphere,
and that intersects the surface in a circle (shown as a red ellipse, the upper half of which
is dashed). Suppose the flare angle q of the cone is such that the area A of the spherical
segment within the circle is equal to one meter squared (A = 1 m2). Then the flare angle
of the cone is equal to 1 steradian (q = 1 sr). The total surface area of the sphere is, in
this case, 12.5664 square meters (4 pi times the square of the radius). Based on the
foregoing example, the geometry of which is independent of scale, it can be said that a
solid angle of 1 sr encompasses about 1/12.5664, or 7.9577 percent, of the space
surrounding a point. The number of steradians in a given solid angle can be determined
by dividing the area on the surface of a sphere lying within the intersection of that solid
angle with the surface of the sphere (when the focus of the solid angle is located at the
center of the sphere) by the square of the radius of the sphere.

STIRLING’S APPROXIMATION: This approximation states that

𝑛 𝑛
𝑛! → 2𝜋𝑛
𝑒
STOCHASTIC: A stochastic variable is the same as a random variable. Generally,
stochastic (pronounced stow-KAS-tik, from the Greek stochastikos, or "skilled at
aiming," since stochos is a target) describes an approach to anything that is based on
probability. In mathematics, a stochastic approach is one in which values are obtained
from a corresponding sequence of jointly distributed random variables. Classic
examples of the stochastic process are guessing the length of a queue at a stated time
given the random distribution over time of a number of people or objects entering and
leaving the queue and guessing the amount of water in a reservoir based on the random
distribution of rainfall and water usage.

Stochastic optimization is the process of maximizing or minimizing the value of a

mathematical or statistical function when one or more of the input parameters is subject to
randomness. The word stochastic means involving chance or probability.

Stochastic processes are commonly involved in business analytics (BA), sales, service,
manufacturing, finance and communications. Stochastic processes always
involveprobability, such as trying to predict the water level in a reservoir at a certain time
based on random distribution of rainfall and water usage, or estimating the number of
dropped connections in a communications network based on randomly variable traffic but
constant available bandwidth. In contrast, deterministic processes never involve probability;
outcomes occur (or fail to occur) based on predictable and exact input values.

STOCHASTIC OPTIMIZATION: Stochastic optimization lends itself to real-life situations

because many phenomena in the physical world involve uncertainty, imprecision or
randomness. Consider the following example: A computer repair shop wants to order
exactly the right number of spare parts of several different types every month to keep
pace with customer demand. If the shop orders too many parts of any type from the
wholesalers, money will be spent needlessly; if the shop does not order enough parts of
any type, it will lose business when customers go elsewhere for service. Determining
the ideal number of parts of each type to order involves stochastic optimization,
because the number of customers who come in with component failures of various sorts
cannot be precisely predicted. The objective is to maximize the function's output value
(the shop's profit) in the face of numerous random input variables.

STOKE’S THEOREM: Stoke’s theorem states that the circulation round any closed curve
T drawn in a fluid is equal to the surface integral 𝑆 of the normal component of spin
taken over any surface, provided the surface lies wholly in the fluid.
 q. dr = curl q ds
T S

Or T=∫S ω . n ds, where n is the unit normal vector at any point of S.

STRATEGY: For a given player the strategy is given by the set of rules which specify that
which of the available course of action he should make at each play. The strategy may be
of two types:

(a) Pure Strategy (b) Mix strategy

(a) PURE STRATEGY : A pure strategy is that in which one player knows what the
other player is going to do. In this case , the player always choose a particular course of
action.
(b) MIXED STRATEGY: When a player does not know exactly what the other player is
going to do, a probabilistic situation is obtained and such type of strategy is known as
mixed strategy . Mathematically , let 𝑥𝑖 be the probability to choose the 𝑖 𝑡𝑕 activity, then
we define the set
𝑋 = 𝑥1 , 𝑥2 , 𝑥3 , … , 𝑥𝑛
𝑠. 𝑡. 𝑥1 +𝑥2 + 𝑥3 , … , +𝑥𝑛 = 1

𝑎𝑛𝑑 𝑥1 , 𝑥2 , 𝑥3 , … , 𝑥𝑛 ≥ 0.

STREAK LINES: A streak line is a lien on which lie all those fluid elements that at some
earlier instant passed through a particular point in space. It is a lien making the position
of a set of fluid particles that had passed through a fixed point in the flow field. A streak
line connects the locations at one instead of particles moving with the fluid which
passed through a particular point. The coloured dye injected in a fluid stream or the
smoke particles rising from the cigarette end exhibits streak lines. When the flow is
steady, streak line coincides with path lines and stream lines.

STREAM FUNCTION: A stream function is a useful — and frequently used — way of

describing a two-dimensional (2-D) incompressible velocity field. The velocity field,
when written in terms of a stream function, automatically satisfies the incompressibility
constraint 𝛻. 𝑢 = 0. For 2-D incompressible flow, we have

𝛻. 𝑢 = 𝜕𝑢/ 𝜕𝑥 + 𝜕𝑣 /𝜕𝑦 = 0
STREAMLINE FLOW: A streamline is a continuous line of flow drawn in the fluid so that
the tangent at every point is the direction of the fluid velocity at the point at given
instant. The component of velocity at right angles to the streamline is always zero. This
shows that there is no flow across the streamline. Thus a solid boundary is also a
streamline.

STREAM SURFACE: A steam surface is a surface made by the stream lines passing
through an arbitrary line in the fluid region at any instant of time.

STREAM TUBE: A stream tube is obtained by drawing the stream lines through every
point of a closed curve in the fluid. Stream surface and stream tubes remain unchanged
in steady motion.

STRENGTH OF THE VORTEX: Let k be the circulation around any closed curve C. Then

k= q dr
c

Consider the closed curve C encloses the vortex tube and lies on its wall, then by Stoke’s
theorem, we have

k= q dr = n curl q ds = ω ∙ n dS = 𝑐𝑜𝑛𝑠𝑡.
c c c

It follows that the circulation round any closed curve enclosing a vortex tube is constant
all along the tube. This constant is also known the strength of the vortex tube. Let 𝜔
denotes the angular velocity and 𝐴 be the cross- section of the vortex tube (assumed
small) then the circulation round this section is 2ωA, which is constant for all section.
This product is called the strength of the vortex.

STRESS TENSOR:- The state of stress at any point of medium is completely characterized
by the specification of nine quantities, called the components of stress tensor.

Let 𝑃 𝑥 be any point in the medium and let → be the stress vector acting on an element
𝑇

of surface at P, with the normal v. Draw through three planar elements parallel to the
coordinate planes and pass the fourth plane ABC normal to v and at a small distance h
from P.
Let the three planar elements which are normal to 𝑥1 , 𝑥2 , 𝑥3 - axes be PBC, PAC,PAB
respectively .

Let →,→ 𝑎𝑛𝑑 → be the stresses acting on the faces 𝒫ℬ𝒞, 𝒫𝒜𝒞 𝑎𝑛𝑑 𝒫𝒜ℬ on the
𝑇 𝑇 𝑇

tetrahedron PABC respectively. Thus → is the stress vector acing on a planar surface
𝑇

element normal to the 𝑥𝑖 - axis .

Resolving the stress vector → into components along the coordinate axes, we have
𝑇

→ = 𝑒𝑗
𝑇 𝑇𝑗

Let us write 𝑇𝑗 = 𝜏𝑖𝑗 𝑠𝑜 𝑡𝑕𝑎𝑡

→ = 𝑒𝑗 𝜏𝑖𝑗 = 𝑒1 𝜏𝑖1 + 𝑒2 𝜏𝑖2 + 𝑒3 𝜏𝑖3

𝑇

𝐼𝑓 𝑖 = 1,2, 3, 𝑤𝑒 𝑕𝑎𝑣𝑒

→ = 𝑒1 𝜏11 + 𝑒2 𝜏12 + 𝑒3 𝜏13 ,

𝑇

→ = 𝑒1 𝜏21 + 𝑒2 𝜏22 + 𝑒3 𝜏23 ,

𝑇

→ = 𝑒1 𝜏31 + 𝑒2 𝜏32 + 𝑒3 𝜏33

𝑇

The nine scalar quantities 𝒯𝑖𝑗 are the components of a stress –tensor . Hence the state of
a stress at a point is completely determined if the nine components of stress tensor at
that point are known.

STRESS VECTOR: Consider an element ∆ 𝜍 of a surface which lies either in the interior or
on the boundary of the medium. Let the force acting on the element ∆ 𝜍 𝑏𝑒 → ∆ 𝜍 .
𝑇

Because of the assumed continuity of forces , we have

∆𝜍
Lim 𝑇 = → 𝑥1 , 𝑥2 , 𝑥3 ,
∆𝜍 𝑇

Where → represents the surface per unit area of the surface acting at the point
𝑇

𝑥𝑖 𝑖. 𝑒. 𝑥1 , 𝑥2 , 𝑥3 and is called the stress- vector

𝑠𝑢𝑟𝑓𝑎𝑐𝑒 𝑓𝑜𝑟𝑐𝑒
Therefore, stress vector → =
𝑇 𝑎𝑟𝑒𝑎
STRICT CONVERSE DUALITY THEOREM: Let X 0 be an open set in Rn and let 𝜃 and 𝑔 be
differentiable and convex on X 0 , let the (primal) minimization problem have a solution
x, and let g satisfy any one of the six constraint qualifications of Kuhn- Tucker stationary
point necessary optimality theorem. If x, u is a solution of the dual (maximization)
problem and if Ψ x, u is strictly convex at x, then x = x that is, x solve MP i.e., (the
primal) minimization problem and θ x = Ψ(x, u).

STRICTLY QUASICONCAVE FUNCTION: A numerical function θ defined on a set T ⊂ Rn is

said to be strictly quasiconcave at x ∇ T (with respect to T) if for each x ∇ T such that
θ x < 𝜃(𝑥), the function θ assumes a higher values then θ(x) on each point in the
intersection of the open line segment (x, x)] and 𝑇 , or equivalently

x ∇ T, θ x < 𝜃 x ,0 < 𝜆 < 1 & 1 − 𝜆 x + λx ∇ T ⇒ θ(x) < 𝜃 1 − 𝜆 x + λx

𝜃 is said to be strictly quasiconcave on 𝑇 if it is strictly quasiconvex at each x ∇ T.

STRICTLY QUASICONVEX FUNCTION: A numerical function θ defined on a set T ⊂ Rn is

said to be quasiconcave at x ∇ T (with respect to T) if for each x ∇ T such that
θ x < 𝜃(x), the function θ assumes a lower values then θ(x) on each point in the
intersection of the open line segment (x, x) and 𝑇 , or equivalently

x ∇ T, θ x < 𝜃 x , 0 < 𝜆 < 1and 1 − 𝜆 x + λx ∇ T ⇒ 𝜃 1 − 𝜆 x + λx < 𝜃(x)

𝜃 is said to be strictly quasiconvex on 𝑇 if it is strictly quasiconvex at each x ∇ T.

STRONG CONVERGENCE: The convergence with respect to the norm is called strong
convergence. Thus the sequence 𝑓𝑛 converges strongly to 𝑓 if

𝑓𝑛 − 𝑓 ⟶ 0, 𝑎𝑠 𝑛 ⟶ ∞

STURM LIOUVILLE EQUATIONS: Various important orthogonal sets of functions arise in

solutions of second-order differential equations of the form

R(x)y ′ + Q x + λP(x) y = 0 ……….. 1

On some interval a ≤ x ≤ b satisfying conditions of the form

a a1 y + a2 y ′ = 0 at x = a
………. 2
and b b1 y + b2 y ′ = 0 at x = b
here λ is a real parameter and a1 , a2 , b1 , b2 , are given real constants at least one in each
conditions (2) being different from zero.

The equation (1) is known as the Sturm- Liouville equation.

SUBALGEBRA: A subspace M of an algebra A is linear subalgebra if

𝑓, 𝑔 𝜖 𝑀 ⟹ 𝑓𝑔 ∇ 𝑴, 𝑓𝑜𝑟 𝑎𝑙𝑙 𝑓, 𝑔 ∇ 𝑀.
SUBGRAPH: Let (𝑉, 𝐸) and(𝑉0 , 𝐸0 ) be graphs. The graph (𝑉0 , 𝐸0 ) is said to be a
subgraph of (𝑉, 𝐸) if and only if 𝑉0 ⊂ 𝑉 and 𝐸0 ⊂ 𝐸 (i.e., if and only if the vertices and
edges of (𝑉0 , 𝐸0 ) are all vertices and edges of (𝑉, 𝐸)).
SUBMANIFOLDS IN 𝑹𝒌 : Let 𝑆 be a manifold in 𝑹𝒌 . A submanifold of 𝑆 is a subset 𝑇 ⊂ 𝑆,
which is also a manifold in 𝑹𝒌 . Let 𝑇 be a submanifold of 𝑆 ⊂ 𝑹𝒌 . Then 𝑑𝑖𝑚𝑇 ≤ 𝑑𝑖𝑚𝑆
and 𝑇𝑝 𝑇 ⊂ 𝑇𝑝 𝑆 for all 𝑝 ∇ 𝑇. Moreover, the inclusion map 𝑖: 𝑇 → 𝑆 is smooth, and its
differential 𝑑𝑖𝑝 at 𝑝 ∇ 𝑇 is the inclusion map 𝑇𝑝 𝑇 → 𝑇𝑝 𝑆.
SUBMATRICES OF A MATRIX: Any matrix obtained by omitting some rows and columns
from a given 𝑚 × 𝑛 matrix 𝐴 is called a submatrix of 𝐴.

The matrix 𝐴 itself is a sub-matrix of 𝐴 as it can be obtained from 𝐴 by omitting no rows

or columns.

A square submatrix of a square matrix 𝐴 is called a principle submatrix, if its diagonal

elements are also the diagonal elements of the matrix 𝐴. Principle submatrices are
obtained only by omitting corresponding rows and columns.

1 2 3 9
1 2 3
Example. The matrix is submatrix of the matrix 𝐴 = 7 11 6 5 as it can
0 2 1
0 2 1 8
be obtained from 𝐴 by omitting the second row and the fourth column.

Equality of two matrics:

Two matrices 𝐴 = 𝑎𝑖𝑗 and 𝐵 = 𝑏𝑖𝑗 are said to be equal if

(iii) They are of the same size and

(iv) The elements in the corresponding places of the two matrices are the same
i.e., 𝑎𝑖𝑗 = 𝑏𝑖𝑗 for each pair of subscripets 𝑖 and 𝑗.
If two matrices 𝐴 and 𝐵 are equal. We write 𝐴 = 𝐵. If two matrices 𝐴 and 𝐵 are not
equal, we write 𝐴 ≠ 𝐵. If two matrices are not of the same size, they cannot be equal.

SUBMERSION: A smooth map 𝑓 ∶ 𝑀 → 𝑁 is called submersion at 𝑃 if the linear

transformation
𝑇𝑓 ∶ 𝑇𝑃(𝑀) → 𝑇𝑓(𝑃)(𝑀)
is surjective.
SUBMODULE: If (𝑀, 𝑅, +, . ) is a module over a ring 𝑅 we say that a subset 𝑁 of 𝑀 is a
submodule if 𝑁 is a subgroup of (𝑀, +) and 𝑟𝑛 ∇ 𝑁 whenever 𝑟 ∇ 𝑅 and 𝑛 ∇ 𝑁. Let 𝑀
be a module over some unital commutative ring 𝑅. Given any subset 𝑋 of 𝑀, the
submodule of 𝑀 generated by the set 𝑋 is defined to be the intersection of all
submodules of 𝑀 that contain the set 𝑋. It is therefore the smallest submodule of 𝑀 that
contains the set 𝑋. An 𝑅 module 𝑀 is said to be finitely-generated if it is generated by
some finite subset of itself.

SUBRINGS: A subsets 𝑆 of a ring 𝐴 is called a subring of 𝐴 if a ring structure is given on 𝑆

and the canonical injection 𝑆 → 𝐴 is a ring homomorphism. Thus the ring operations of
𝑆 are the restrictions of those of 𝐴. If we deal only with unitary rings and unitary
homomorphisms, then a subring 𝑆necessarily contains the unity element of 𝐴. The
smallest subring containing a subset 𝑇 of a ring 𝐴 is called the subring generated by 𝑇.
The set of elements that commute with every element of 𝑇 forms a subring and is called
the commuter (or centralizer) of 𝑇. In particular, the commuter of 𝐴 itself is called the
center of 𝐴.

SUBSET: Set 𝐵 is a subset of set 𝐴 if every element contained in 𝐵 is also contained in 𝐴.

SUBSETHOOD IN FUZZY SETS: Let 𝐴 and 𝐵 are fuzzy subsets of a classical set 𝑋. We say
that 𝐴 is a subset of 𝐵 if 𝐴(𝑡) ≤ 𝐵(𝑡), ∀𝑡 ∇ 𝑋.
SUBSPACE OF A NORMED SPACE: By a subspace of a normed space (or inner product
space) we mean a linear subspace with the same norm (inner product respectively). We
write 𝑋 ⊂ 𝑌 or 𝑋 ⊆ 𝑌.

 𝐶𝑏 (𝑋) ⊂ 𝑙∞ (𝑋) where 𝑋 is a metric space.

 Any closed subspace of a Banach/Hilbert space is complete, hence also a
Banach/Hilbert space.
 Any complete subspace is closed.
  The closure of subspace is again a subspace.

The space 𝐶[𝑎, 𝑏] is incomplete for any 𝑎 < 𝑏 if equipped by the inner product and the
1
𝑏 2 𝑏 2 2
corresponding norm: ⌌𝑓, 𝑔⌍ = ∫𝑎 𝑓 𝑡 ḡ t dt, 𝑓 = ∫𝑎 𝑓(𝑡) 𝑑𝑡

SUB-SPACE OF AN 𝒏-VECTOR SPACE 𝑽𝒏 : A non-empty set 𝑆, of vectors of 𝑉𝑛 is called a

vector subspace of 𝑉𝑛 , if 𝑎 + 𝑏 belong to 𝑆 whenever 𝑎, 𝑏 belong to 𝑆 and 𝑘𝑎 belongs to 𝑆
whenever a belongs to 𝑆, where 𝑘 is any scalar.

It is important to note that every sub-space of 𝑉𝑛 contains the zero vector, being the
scalar product of any vector with the scalar zero.

SUBSPACES: Let V be a vector space over the field K. Certain subsets of V have the nice
property of being closed under addition and scalar multiplication; i.e., adding or taking
scalar multiples of vector in the subset gives vectors which are again in the subset. We
call such a subset a subspace:

A subspace of V is a non-empty subset W ⊆ V such that

(i) W is closed under addition: 𝒖, 𝒗 𝜖 𝑊 ⟹ 𝒖 + 𝒗 𝜖 𝑊

(ii) W is closed under scalar multiplication: 𝒗 𝜖 𝑊, 𝛼 𝜖 𝐾 ⟹ 𝛼 𝒗 𝜖 𝑊
These two conditions can be replaced with a single condition

𝒖, 𝒗 ϵ W , 𝛼, 𝛽 𝜖𝐾 ⟹ 𝛼𝐮 + 𝛽 𝒗 𝜖 𝑊

A subspace W is itself a vector space over K under the operations of vector addition and
scalar multiplication in V. Notice that all vector space axioms of W hold automatically.

Example: The subset of ℝ2 given by W= 𝛼, 𝛽 𝜖 ℝ2 : 𝛽 = 2𝛼 i.e the subset consisting of

all row vectors whose second entry is twice their first entry, is a subspace of ℝ2 .Adding
twovectors of this form always gives another vector of this form and multiplying a
vector of this form bt a scalar always gives another vector of this form.

For any Vector space V, V is always a subspace of itself. Subspaces other than V are
sometimes called proper subspaces. We also always have a subspace 0 consisting of
the zero vector alone.This is called the trivial subspace, and its dimension is 0, because
it has no linearly independent set of vectors at all.
Intersection of two subspaces gives third subspace i.e. If 𝑊1 𝑎𝑛𝑑 𝑊2 are subspaces of V
then so is 𝑊1 ∩ 𝑊2 .

It is not necessarily true that 𝑊1 ∪ 𝑊2 is a subspace, as the following example shows:

Example: Let V=ℝ2 , let 𝑊1 = 𝛼, 0 : 𝛼 𝜖ℝ 𝑎𝑛𝑑 𝑊2 = 0, 𝛼 : 𝛼 𝜖ℝ .Then 𝑊1 , 𝑊2

Are subspaces of V but 𝑊1 ∪ 𝑊2 is not a subspace .because (1,0),(0,1) 𝜖𝑊1 ∪ 𝑊2 but

(1,0)+(0,1)=(1,1) ∈ 𝑊1 ∪ 𝑊2 .

Note that any subspace of V that contains 𝑊1 𝑎𝑛𝑑 𝑊2 has to contain all vectors of the
form u+v for u 𝜖 𝑊1 , 𝑣 𝜖 𝑊2

Let 𝑊1 , 𝑊2 be two subspaces of the Vector space V.Then 𝑊1 + 𝑊2 is defined to be the

set of vectors v 𝜖 𝑉 such that v= 𝑤1 + 𝑤2 for some 𝑤1 𝜖 𝑊1 , 𝑤2 𝜖 𝑊2

Or 𝑊1 + 𝑊2 = 𝑤1 + 𝑤2 : 𝑤1 𝜖 𝑊1 , 𝑤2 𝜖 𝑊2

If 𝑊1 , 𝑊2 are subspaces of the Vector space V, then 𝑠𝑜 𝑖𝑠 𝑊1 + 𝑊2 .Infact it is the

smallest subspace that contains both 𝑊1 𝑎𝑛𝑑 𝑊2 .

Let V be a finite-dimensional vector space and let 𝑊1 , 𝑊2 are subspaces of V. Then

𝑑𝑖𝑚 (𝑊1 + 𝑊2 ) = 𝑑𝑖𝑚𝑊1 + 𝑑𝑖𝑚𝑊2 − 𝑑𝑖𝑚⁡

(𝑊1 ∩ 𝑊2 )

Let 𝑣1 , 𝑣2 , ⋯ ⋯ , 𝑣𝑛 be vectors in the vector space V.Then the set of all linear
combinations 𝛼1 𝑣1 + 𝛼2 𝑣2 + ⋯ ⋯ + 𝛼𝑛 𝑣𝑛 of 𝑣1 , 𝑣2 , ⋯ ⋯ , 𝑣𝑛 forms a subspace of V.

SUBTRACTION OF TWO MATRICES: If A and B are two 𝑚 × 𝑛 matrices, then we define

A − B = A + (−B).

Thus the difference A − B is obtained by subtracting from each elements of A the

corresponding elements of B.

SUCCESSOR OF A SET: Given a set 𝑆, the successor of 𝑆 is the set 𝑆 ′ = 𝑆 ∪ {𝑆}. One often
denotes the successor of 𝑆 by 𝑆 1 .
SUFFICIENT OPTIMALITY THEOREM (KUHN- TUCKER): Let x ∇ X 0 , let X 0 be open and let
θ and g be differentiable and convex at x. If x, u is a solution of Kuhn- Tucker problem:

∆θ x) + u ∆g(x = 0, g(x) ≤ 0,
 u g x = 0, u ≥ 0,

Then x solves the problem: Minimize 𝜃 x subject to x ∇ X 0 , g(x) ≤ 0.

If x, r0 , r is a solution of Fritz John problem:

r0 ∆θ x) + r ∆g(x = 0, g(x) ≤ 0, r g x = 0,(r0 , r) ≥ 0 and r0 > 0, then x is a solution of

the problem: Minimize 𝜃 x subject to x ∇ X 0 , g(x) ≤ 0.

SUM AND SCALAR PRODUCT OF MEASURE: Let µ1 and µ2 be measures on a same ς-

algebra. Define µ1 + µ2 and 𝜆µ1 , 𝜆 > 0 by (µ1 + µ2 )(𝐴) = µ1 (𝐴) + µ2 (𝐴) and
(𝜆µ1 )(𝐴) = 𝜆(µ1 (𝐴)). Then µ1 +µ2 and 𝜆µ1 are measures on the same 𝜍 −algebra as
well.

SUMMABLE FUNCTION: A simple function with disjoint 𝐴𝑘 is called summable if the

∞ ∞
series 𝑘=1 𝑡𝑘 𝜇(𝐴𝑘 ) converges, where 𝑓 has the unique representation 𝑓 = 𝑘=1 𝑡𝑘 𝜒𝐴𝑘

SUPERSET: Given two sets 𝐴 and 𝐵, 𝐴 is a superset of 𝐵 if every element in 𝐵 is also in 𝐴.

We denote this relation as 𝐴 ⊇ 𝐵. This is equivalent to saying that 𝐵 is a subset of 𝐴,
that is 𝐴 ⊇ 𝐵 ⇒ 𝐵 ⊆ 𝐴. Similar rules that hold for ⊆ also hold for ⊇. If 𝑋 ⊆ 𝑌 and
𝑌 ⊆ 𝑋, then 𝑋 = 𝑌 . Every set is a superset of itself, and every set is a superset of the
empty set. 𝐴 is a proper superset of 𝐵 if 𝐴 ⊇ 𝐵 and 𝐴 ≠ 𝐵. This relation is often
denoted as 𝐴 ⊃ 𝐵.
SUPPLEMENTARY: Two angles are supplementary if the sum of their measures is 180°.
For example, two angles measuring 135° and 45° form a pair of supplementary angles.
SUPPORT OF A FUZZY SET: The support of a fuzzy set 𝐴 in the universe of discourse 𝑈 is
a crisp set that contains all the elements of 𝑈 that have nonzero membership values in
𝑨, that is,
𝑠𝑢𝑝𝑝(𝐴) = {𝒂 ∇ 𝑈; 𝝁𝑨 (𝑥) > 0}
SURFACE: A surface is said to be locus of a point whose Cartesian coordinates (𝑥, 𝑦, 𝑧)
are functions of two independent, say (𝑢, 𝑣).
i.e., 𝑥 = 𝑓 𝑢, 𝑣 ; 𝑦 = 𝑔 𝑢, 𝑣 ; 𝑧 = 𝑕(𝑢, 𝑣)
above equations are called parametric or freedom equations of a surface, further the
parameters 𝑢, 𝑣 assume real values and vary in some region say 𝐷. We may define a
surface vectorially as the locus of a point whose position vector 𝑟 is expressed in terms
of parameters i.e., 𝑟 = 𝑟(𝑢, 𝑣) represents a surface.
The above representation of the surface is due to Gauss and is termed as Gaussian form
of the surface. The parameters 𝑢 and 𝑣 are called curvilinear coordinates or surface co-
ordinates of the current point on the surface. The pair of parameters 𝑢 and 𝑣 represents
the points (𝑢, 𝑣). What is called a surface or a curved surface is usually a 2-dimensional
topological manifold, that is, a topological space that satisfies the second countability
axiom and of which every point has a neighborhood homeomorphic to the interior of a
circular disk in a 2-dimensional Euclidean space. A compact surface is called a closed
surface, and a noncompact surface is called an open surface. A closed surface is
decomposable into a finite number of 2-simplexes and so can be interpreted as a
combinatorial manifold. A 2-dimensional topological manifold having a boundary is
called a surface with boundary.
SURFACE AREA: The surface area of a solid is a measure of how much area the solid
would have if you could somehow break it apart and flatten it out. For example, a cube
with edge a units long has six faces, each with area 𝑎2 . The surface area of the cube is
the sum of the areas of these six faces, or 6𝑎2 . The surface area of any polyhedron can be
found by adding together the areas of all the faces. The surface areas of curved solids
are harder to find, but they can often be found with calculus.
SURFACE AREA, FIGURE OF REVOLUTION: Suppose the curve 𝑦 = 𝑓(𝑥) is rotated about
the 𝑥 −axis between the lines 𝑥 = 𝑎 and 𝑥 = 𝑏. The surface area of this figure can be
found with integration. Let dA represent the surface area of a small frustum cut from
this figure. The surface area of the frustum is 𝑑𝐴 = 2𝑝𝑦𝑑𝑠 where 𝑦 is the average
radius of the frustum, and ds is the slant height. 𝑑𝑆 is given by the formula

2
𝑑𝑦
𝑑𝑆 = 1+ 𝑑𝑥 .
𝑑𝑥
Then the total surface area is given by this integral:
𝑏
2
𝑑𝑦
2𝜋𝑦 1 + 𝑑𝑥.
𝑑𝑥
𝑎

SURFACE FORCES: The forces which act on the surface of the volume elements are called
surface forces.

SURFACE INTEGRAL: Let E be a three-dimensional vector field, and let S be a surface.

Consider a small square on this surface. Create vector dS whose magnitude is equal to
the area of the small square, and whose direction is oriented to point outward along the
surface. Calculate the dot product 𝑬. 𝒅𝑺 and then integrate this dot product over the
entire surface. The result is the surface integral of the field E along this surface:

𝑬. 𝒅𝑺

SURFACE OF REVOLUTION: A surface generated by the revolution of a plane curve

about axis in its plane is called a surface of revolution.

SURJECTIVE FUNCTION: We say that the function f is surjective if, given any element 𝑏 of
𝐵, there exists some element 𝑎 of 𝐴 such that 𝑓(𝑎) = 𝑏. Thus a function is surjective if
every element of the codomain is the image of some element of the domain. We say that
the function f is bijective if it both injective and surjective.

SURPLUS VARIABLES: If the given constraints in LPP have a sign ≥ , then to turn the
inequations as equality, a positive quantity is to be subtracted from LHS. The positive
quantities (or variables), that are to be subtracted from LHS to form the equations are
called surplus variables.

SUSLIN’S HYPOTHESIS (SH): Every dense, linear, order complete set without end points,
having at most 𝑤 disjoint intervals, is order isomorphic to the continuum of real
numbers.
SYLOW p-SUBGROUPS: Let 𝐺 be a finite group and 𝑝 be a prime that divides |𝐺|. We can
then write |𝐺| = 𝑝𝑘 𝑚 for some positive integer 𝑘 so that 𝑝 does not divide 𝑚. Any
subgroup of 𝐻 whose order is 𝑝𝑘 is called a Sylow 𝑝-subgroup or simply 𝑝 subgroup.
First Sylow theorem states that any group with order 𝑝𝑘 𝑚 has a Sylow 𝑝-subgroup.
SYMMETRIC: (1) Two points A and B are symmetric with respect to a third point (called
the center of symmetry) if the third point is the midpoint of the segment connecting the
first two points.
(2) Two points A and B are symmetric with respect to a line (called the axis of
symmetry) if the line is the perpendicular bisector of the segment AB.
SYMMETRIC CHANNEL: A symmetric channel is a memoryless communication channel
1
for which there exists a 𝑝 < 2 such that for every 𝑖, 𝑗 ∇ {0, 1}𝑛 with 𝑖 ≠ 𝑗 it holds that
𝑞

𝑃[𝑎𝑗 𝑟𝑒𝑐𝑒𝑖𝑣𝑒𝑑 | 𝑎𝑖 𝑤𝑎𝑠 𝑠𝑒𝑛𝑡] = 𝑝

𝑗 =1,𝑗 ≠ 𝑖
 th
SYMMETRIC MATRIX: A square matrix A = aij is said to be symmetric if its i, j
element i.e., if aij = aij for all i. j.

a h g p
h b f g 1 i −2i
2 4
For example, . i −2 4 , are symmetric matrices
g f c r 4 3
−2i 4 3
p q r s

A necessary and sufficient condition for a matrix A to be symmetric is that A and A’ are
equal.

SYMMETRIC POLYNOMIAL: A polynomial 𝑓 ∇ 𝑅[𝑥1 , . . . , 𝑥𝑛 ] in 𝑛 variables with

coefficients in a ring 𝑅 is symmetric if 𝛼(𝑓) = 𝑓 for every permutation 𝛼 of the set
𝑥1 , . . . , 𝑥𝑛 . Every symmetric polynomial can be written as a polynomial expression in
the elementary symmetric polynomials.
SYMMETRIC PROPERTY OF EQUALITY: The symmetric property of equality states that, if
𝑎 = 𝑏 , then 𝑏 = 𝑎 . That means that you can reverse the two sides of an equation
whenever you want to.
SYMMETRIC PROPERTY OF HYPERGEOMETRIC FUNCTION: Hypergeometric function
does not change if the parameters α and β are interchanged, keeping γ fixed.

SYMMETRIC TENSOR: A tenser 𝐴𝑖𝑗 is said to be symmetric if 𝐴𝑖𝑗 = 𝐴𝑗𝑖 . Similarly a tensor
𝐴𝑖𝑗𝑘 is symmetric in the suffixes 𝑗 and 𝑘 if 𝐴𝑖𝑗𝑘 = 𝐴𝑗𝑘𝑖 .

SYMMETRY PRINCIPLE: If 𝑓 + and 𝑓 − are holomorphic functions in 𝛺 + and 𝛺 −

respectively, that extend continuously to 𝐼 and
𝑓 +(𝑥) = 𝑓 −(𝑥) ∀ 𝑥 ∇ 𝐼,
then the function f defined on 𝛺 by
𝑓 +(𝑥)∀ 𝑥 ∇ 𝛺 +
𝑓 𝑧 = 𝑓 +(𝑥) = 𝑓 −(𝑥) ∀ 𝑥 ∇ 𝐼
𝑓 −(𝑥) ∀ 𝑥 ∇ 𝛺 −
is holomorphic on all of 𝛺.
SYNDROME DECODING: Let 𝐶 be a linear [𝑛, 𝑘, 𝑑]-code over 𝐹𝑞 and let 𝐻 be a parity-
check matrix for 𝐶. Then for every 𝑤 ∇ 𝐹𝑞𝑛 the syndrome of 𝑤 determined by 𝐻 is
defined to be the word 𝑆(𝑤) = 𝑤𝐻 𝑇 ∇ 𝐹𝑞𝑛−𝑘
SYSTEM OF INEQUALITIES: A system of inequalities is a group of inequalities that are all
to be true simultaneously. For example, this system of three inequalities
𝑥> 3
𝑦> 4
𝑥 + 𝑦 < 15
defines a set of values for x and y that will make all of the inequalities true.
TANGENT BUNDLE: The tangent bundle of 𝑀 is the union 𝑇𝑀 = ⋃𝑝∇𝑀 𝑇𝑝 𝑀 of all
tangent vectors at all points. For a given element 𝑋 ∇ 𝑇𝑀 the point 𝑝 ∇ 𝑀 for which
𝑋 ∇ 𝑇𝑝 𝑀 is called the base point. The map 𝜋: 𝑇𝑀 → 𝑀 which assigns 𝑝 to 𝑋, is called
the projection. Notice that the union we take is disjoint, that is, there is no overlap
between 𝑇𝑝 1 𝑀 and 𝑇𝑝 2 𝑀 if 𝑝1 ≠ 𝑝2 . Formally, an element in 𝑇𝑀 is actually a pair (𝑝, 𝑋)
where 𝑝 ∇ 𝑀 and 𝑋 ∇ 𝑇𝑝 𝑀, and π maps (𝑝, 𝑋) to the first member of the pair.
Notationally it is too cumbersome to denote elements in this fashion, and hence the base
point 𝑝 is suppressed.
TANGENT CIRCLES: Two circles are tangent if they touch at just one point.

TANGENT SPACE OF A PARAMETRIZED MANIFOLD: Let 𝜍: 𝑈 → 𝑅 𝑛 be a parametrized

manifold, where 𝑈 ⊂ 𝑅 𝑚 is open. The tangent space 𝑇𝑥 ∘ 𝜍 of 𝜍 at 𝑥0 ∇ 𝑈 is the linear
subspace of 𝑅 𝑛 spanned by the columns of the 𝑛 × 𝑚 matrix 𝐷𝜍(𝑥0 ).
TAUTOLOGY: A tautology is a sentence that is necessarily true because of its logical
structure, regardless of the facts.
TAYLOR’S THEOREM COMPLEX ANALYSIS): If 𝑎 function 𝑓(𝑧) is analytic within
𝑎 𝑐𝑖𝑟𝑐𝑙𝑒 𝐶 with its centre 𝑧 = 𝑎 radius 𝑅, then at every point 𝑧 inside 𝐶.
 ∞ ∞
(𝑛)
(𝑧 − 𝑎)𝑛
𝑓 𝑧 = 𝑓 𝑎 𝑜𝑟 𝑓 𝑧 = 𝑎𝑛 (𝑧 − 𝑎)𝑛 ,
𝑛!
𝑛=0 𝑛=0

𝑓 𝑛 (𝑎)
Where 𝑎𝑛 = 𝑛!

TAYLOR’S THEOREM WITH CAUCHY’S FORM OF REMAINDER: If 𝑓(𝑥) is a single valued

function of 𝑥 such that
 all the derivatives of 𝑓(𝑥) upto (𝑛 − 1)th order are continuous in the closed
interval [𝑎, 𝑎 + 𝑕]
 𝑓 (𝑛 ) (𝑥) exists in (𝑎, 𝑎 + 𝑕)
then
𝑕2 ′′ 𝑕𝑛 −1
𝑓 𝑎 + 𝑕 = 𝑓 𝑎 + 𝑕𝑓’ 𝑎 + 𝑓 𝑎 + − − − + 𝑓 𝑛 −1 𝑎
2! (𝑛 − 1)!
𝑕𝑛 𝑛−1 𝑛
+ 1−𝜃 𝑓 𝑎 + 𝜃𝑕
(𝑛 − 1)!
where 𝜃 ∇ 0, 1
TAYLOR’S THEOREM WITH LAGRANGE’S FORM OF REMAINDER: If 𝑓(𝑥) is a single
valued function of 𝑥 such that
 all the derivatives of 𝑓(𝑥) upto (𝑛 − 1)th order are continuous in the closed
interval [𝑎, 𝑎 + 𝑕]
 𝑓 (𝑛 ) (𝑥) exists in (𝑎, 𝑎 + 𝑕)
then
𝑕2 ′′ 𝑕𝑛 −1 𝑛−1
𝑕𝑛 𝑛
𝑓 𝑎 + 𝑕 = 𝑓 𝑎 + 𝑕𝑓’ 𝑎 + 𝑓 𝑎 + − − − + 𝑓 𝑎 + 𝑓 𝑎 + 𝜃𝑕
2! (𝑛 − 1)! 𝑛!
where 𝜃 ∇ 0, 1
t-DISTRIBUTION: The t-distribution refers to a family of continuous random variables
that play an important part in statistical estimation theory. A specific t-distribution is
characterized by a parameter known as the degrees of freedom. The density function for
the t-distribution is bell-shaped and centered at 0, similar to the standard normal
distribution. As the degrees of freedom increase, the t-distribution density function
approaches the standard
normal density function.
TENSOR: The physical quantities which have more than one direction are represented
by the mathematical entities, are called tensors. Scalars and vectors are special case of
tensors. A tensor is a type of linear function with multiple indices.
A tensor of rank zero can be represented as a scalar; a tensor of rank one can be
represented as a vector; and a tensor of rank two can be represented as a matrix. In
three-dimensional space the components of a tensor of rank n form a multidimensional
array with 3𝑛 numbers that need to be specified.
A specific tensor is applied to an array of specific dimensions, resulting in an another
array of specified dimensions, similar to the way that an ordinary function is applied to
a number resulting in another number.
TENSORS OF HIGHER RANK: Let 𝐴𝑖𝑗 𝑖, 𝑗 = 1,2,3, … , 𝑛) be 𝑛2 functions of co-ordinates
𝑥1 , 𝑥 2 , … 𝑥 𝑛 and let these transform to 𝐴𝑖𝑗 in another co-ordinate system 𝑥′1 , 𝑥′2 , … 𝑥′𝑛
according to the rules

𝜕𝑥 ′𝑖 𝜕𝑥 ′𝑗
𝐴′𝑖𝑗 = 𝐴𝑎𝛽 .
𝜕𝑥 𝑎 𝜕𝑥 𝛽

Then 𝐴𝑖𝑗 are called components of a contravariant tensor of rank two. Similarly if
𝐴𝑖𝑗 𝑖, 𝑗, = 1,2, … 𝑛 be 𝑛2 funcations of ordinates 𝑥1 , 𝑥 2 , … 𝑥 𝑛 and if 𝐴𝑖𝑗 are transformed
to 𝐴′𝑖𝑗 in another co-ordinate system 𝑥′1 , 𝑥′2 , … 𝑥′𝑛 by the rule

𝜕𝑥 𝑎 𝜕𝑥 𝛽
𝐴′𝑖𝑗 = 𝐴𝑎𝛽
𝜕𝑥 ′𝑖 𝜕𝑥 ′𝑗

Then 𝐴𝑖𝑗 are said to be covariant tensor of rank two.

Finally, if the 𝑛2 funcations 𝐴𝑗𝑖 𝑖, 𝑗 = 1,2, … 𝑛 of co-ordinates 𝑥 𝑖 be transformed of 𝐴𝑗′𝑖 in

another co-ordinate system 𝑥′𝑖 according to the rules

𝜕𝑥 ′𝑖 𝜕𝑥 𝛽
𝐴𝑗′𝑖 = 𝐴𝛽𝑎
𝜕𝑥 𝑎 𝜕𝑥 ′𝑗

Then 𝐴𝑗𝑖 are said to be components of mixed tensor of rank two.

TEST STATISTIC: A test statistic is a quantity calculated from observed sample values
that is used to test a null hypothesis. The test statistic is constructed so that it will come
from a known distribution if the null hypothesis is true. Therefore, the null hypothesis is
rejected if it seems implausible that the observed value of the test statistic could have
come from that distribution.
TETRAHEDRON: A tetrahedron is a polyhedron with four faces. Each face is a triangle. In
other words, a tetrahedron is a pyramid with a triangular base. A regular tetrahedron
has all four faces congruent.

1
TETRAHEDRAL NUMBER: An integer of the form 6 𝑛 𝑛 + 1 𝑛 + 2 , where 𝑛 is a positive

integer. This number equals the sum of the first 𝑛 triangular numbers. The first few
tetrahedral numbers are 1, 4, 10 and20.
TFAE: The abbreviation “TFAE” is shorthand for “the following are equivalent”. It is used
before a set of equivalent conditions.
THE IDENTITY THEOREM: Suppose that 𝑓 is meromorphic on a domain 𝐷 in 𝐶. Suppose
that 𝑏 ∇ 𝐶. Then either 𝑓(𝑧) ≡ 𝑏 on 𝐷, or the set 𝐸 = {𝑧 ∇ 𝐷 ∶ 𝑓(𝑧) = 𝑏} has no
limit point in 𝐷.
THEN: The word “THEN” is used as a connective word in logic sentences of the form “p S
q” “IF p, THEN q.” Here is an example: “If a triangle has three equal sides, then it has
three equal angles.”
THE SPACE 𝑩 [𝒂, 𝒃 ]: The space of all those complex valued functions defined on 𝑎, 𝑏
which are such that, for some real number 𝑀, 𝑓(𝑥) ≤ 𝑀, is denoted by 𝐵 𝑎, 𝑏 .

THETA FUNCTIONS: A holomorphic function 𝑓(𝑧) on 𝐶 𝑛 = (𝑅 𝑛 , 𝐽) is called a theta

function if for every 𝑑 ∇ 𝐷, we have 𝑓(𝑧 + 𝑑) = 𝑓(𝑧)𝑒𝑥𝑝(𝑙𝑑 (𝑧) + 𝑐𝑑 ), where 𝑙𝑑 (z) is a
linear form on 𝐶 𝑛 which, as for 𝑐𝑑 , depends 𝑑.
THREE DIMENSIONAL SOURCE AND SINK: If the motion of a liquid consists of
symmetrical outward radial flow in all direction proceeding from a point then the point
is called a three- dimensional source.

Let 4πm is the volume flux per unit time from the source point O then m is called the
strength of the source.
TIME SERIES: A sequence of observations taken over a period of time, usually at equal
intervals. Time series analysis is concerned with identifying the factors that influence
the variation in a time series, perhaps with a view to predicting what will happen in the
future.
TISSOT’S THEOREM DIFFERENTIAL GEOMETRY): Let S, S* be two surface and 𝑓 be a
non conformal mapping of the surface S onto the surface 𝑆 ∗ , given by a differential
homeomorphism regular at each point, there exists at every point P of S a uniquely
determined pair of real orthogonal directions such that the corresponding directions of
S* are also orthogonal.

T-NORM IN FUZZY SETS: Any function 𝑡 ∶ [0, 1] × [0, 1] → [0, 1] that satisfies
following Axioms tl-t4 is called a 𝑡 −norm.
Axiom t1: 𝑡(0,0) = 0; 𝑡(𝑎, 1) = 𝑡(1, 𝑎) = 𝑎 (boundary condition).
Axiom t2: 𝑡(𝑎, 𝑏) = 𝑡(𝑏, 𝑎) (commutativity).
Axiom t3: If 𝑎 ≤ 𝑎′ and 𝑏 ≤ 𝑏′, then 𝑡(𝑎, 𝑏) ≤ 𝑡(𝑎′ , 𝑏′) (non-decreasing condition).
Axiom t4: 𝑡[𝑡(𝑎, 𝑏), 𝑐] = 𝑡[𝑎, 𝑡(𝑏, 𝑐)] (associativity).
TOPOLOGICAL DIVISORS OF ZERO: An element 𝑓 of Banach algebra A is called
topological divisor of zero if there exists a sequence 𝑓𝑛 in A sich that 𝑓𝑛 = 1 and
either 𝑓𝑓𝑛 ⟶ 0 or 𝑓𝑛 𝑓 ⟶ 0. We denote the set of all topological divisors of zero by Z.
TOPOLOGICAL GROUP: We say that (𝐺,×, 𝜏 ) is a topological group if (𝐺,×) is a group
and (𝐺, 𝜏 ) a topological space such that, writing 𝑀(𝑥, 𝑦) = 𝑥 × 𝑦 and 𝐽𝑥 = 𝑥 −1 the
multiplication map 𝑀 ∶ 𝐺 2 → 𝐺 and the inversion map 𝐽 ∶ 𝐺 → 𝐺 are continuous.
TOPOLOGICAL ISOMORPHISM: If (𝐺,×𝐺 , 𝜏𝐺 ) and (𝐻,×𝐻 , 𝜏𝐻 ) are topological groups we
say that 𝜃 ∶ 𝐺 → 𝐻 is an isomorphism if it is a group isomorphism and a topological
homeomorphism.
TOPOLOGICAL MANIFOLD: A Hausdorff space 𝑀 is called a (topological) 𝑛-manifold if
each point of 𝑀 has a neighborhood homeomorphic to an open set in 𝑅 𝑛 . Roughly
speaking, an 𝑛-manifold is locally 𝑅 𝑛 . Sometimes 𝑀 is denoted as 𝑀𝑛 for mentioning the
dimension of 𝑀.
TOPOLOGY: Topology is the mathematical study of how points are connected together. If
an object is stretched or bent, then its geometric shape changes but its topology remains
unchanged. Let 𝑋 be a non-empty set. A set ℑ of subsets of 𝑋 is said to be a topology on
𝑋 if
(i) 𝑋 and the empty set, 𝜑, belong to ℑ ,
(ii) The union of any (finite or infinite) number of sets in ℑ belongs to ℑ ,
(iii) The intersection of any two sets in ℑ belongs to ℑ .
The pair (𝑋; ℑ) is called a topological space.
In other words, a topological space is a set X and a collection of subsets of X called open
sets such that:
(1) The empty set is open,
(2) The whole set is open,
(3) A finite intersection of open sets is open,
(4) An arbitrary union of open sets is open.
A metric space (𝑋, 𝑑) is a topological space: the open sets are any union of balls
𝐵𝑟(𝑥) = {𝑦 ∇ 𝑋 ∶ 𝑑(𝑥, 𝑦) < 𝑟} (for centres 𝑥 ∇ 𝑋, radii 𝑟 > 0).
TOPOLOGY OF R: With the collection of all its open intervals (𝑎, 𝑏) as an open base, 𝑅 is
a topological space (order topology) that satisfies the separation axioms 𝑇2 , 𝑇3 , 𝑇4 . In 𝑅
every (finite or infinite) interval (including 𝑅 itself) is connected, and the set 𝑄 of
rational numbers is dense. A necessary and sufficient condition for a subset 𝐹 of 𝑅 to be
compact is that 𝐹 be bounded and closed Weierstrass’s theorem . In particular, any
finite closed interval is compact. 𝑅 is a locally compact space satisfying the second
countablity axiom. Further, any (finite or infinite) open interval is homeomorphic to 𝑅.

TOROID: A toroid can be formed by rotating a closed curve for a full turn about a line
that is in the same plane as the curve, but does not cross it. The set of all points that the
curve crosses in the course of the rotation forms a toroid.
TORSION: Torsion at point 𝑃of a given curve is the arc rate of the change in the direction
of the binormal at 𝑃. Its magnitude is denoted by 𝜏(𝑇𝑎𝑢).
TORSION-FREE MODULES: A module 𝑀 over an integral domain 𝑅 is said to be
torsionfree if 𝑟𝑚 is non-zero for all non-zero elements 𝑟 of 𝑅 and for all non-zero
elements 𝑚 of 𝑀.
TORSION MODULES: A module 𝑀 over an integral domain 𝑅 is said to be a torsion
module if, given any element 𝑚 of 𝑀, there exists some non-zero element 𝑟 of 𝑅 such
that 𝑟𝑚 = 0𝑀 , where 0𝑀 is the zero element of 𝑀.
TORSION SUBMODULE: Let 𝑀 be a module over an integral domain 𝑅. The torsion
submodule of 𝑀 is the submodule 𝑇 of M defined such that 𝑇 = {𝑚 ∇ 𝑀 ∶ 𝑟𝑚 = 0𝑀
for some non-zero element 𝑟 of 𝑅}, where 0𝑀 denotes the zero element of 𝑀. Thus an
element 𝑚 of 𝑀 belongs to the torsion submodule 𝑇 of 𝑀 if and only if there exists some
non-zero element 𝑟 of 𝑅 for which 𝑟𝑚 = 0𝑀 .
TORUS: A torus is a solid figure formed by rotating a circle about a line in the same
plane as the circle, but not on the circle. A tube of a cycle is an example of a torus.
TOTAL ORDER: A total order is a special case of a partial order. If ≤ is a partial order on
𝐴, then it
satisfies the following three properties:
1. reflexivity: 𝑎 ≤ 𝑎 for all 𝑎 ∇ 𝐴
2. antisymmetry: If 𝑎 ≤ 𝑏 and 𝑏 ≤ 𝑎 for any 𝑎, 𝑏 ∇ 𝐴, then 𝑎 = 𝑏
3. transitivity: If 𝑎 ≤ 𝑏 and 𝑏 ≤ 𝑐 for any 𝑎, 𝑏, 𝑐 ∇ 𝐴, then 𝑎 ≤ 𝑐
The relation ≤ is a total order if it satisfies the above three properties and the following
additional property:
Comparability: For any 𝑎, 𝑏 ∇ 𝐴, either 𝑎 ≤ 𝑏 or 𝑏 ≤ 𝑎.

TOWER OF HANOI: Imagine three poles with a number of discs of different sizes initially
all on one of the poles in decreasing order. The problem is to transfer all the discs to one
of the other poles, moving the discs individually from pole to pole, so that one disc is
never placed above another of smaller
diameter. A version with 8 discs, known as the Tower of Hanoi, was invented by
Edouard Lucas. If 𝑘𝑛 is the number of moves it takes to transfer 𝑛 discs from one pole to
another, then 𝑘1 = 1 and 𝑘𝑛+1 = 2𝑘𝑛 + 1. It can be shown that the solution of this
difference equation is 𝑘𝑛 = 2𝑛 – 1, so the original Tower of Hanoi puzzle takes 255
moves.
TOY THEOREM: A toy theorem is a simplified version of a more general theorem. For
instance, by introducing some simplifying assumptions in a theorem, one obtains a toy
theorem.
TOWER LAW: If K is a subfield of 𝐿 and 𝐿 is a subfield of 𝑀 then [𝑀 ∶ 𝐾] = [𝑀 ∶ 𝐿][𝐿 ∶
𝐾].

TRACE OF A MATRIX: Let 𝐴 be a square matrix of order 𝑛. The sum of the elements of 𝐴
lying along the principal is called the trace of 𝐴. We shall write the trace of 𝐴 as tr 𝐴.
Thus if A = aij , then tr
n×n
 n
A= i=1 a ij = a11 + a22 + ⋯ + ann .

Let A and 𝐵 be two square matrices of order 𝑛 and 𝜆 be a scalar. Then

(1) tr 𝜆𝐴 = 𝜆 tr 𝐴;
(2) tr 𝐴 + 𝐵 = tr 𝐴 + 𝑡𝑟 𝐵;
(3) tr 𝐴𝐵 = tr B𝐴 .

TRACE OF AN ENDOMORPHISM: The trace of an endomorphism 𝜏 of an 𝑛-dimensional

vector space is the trace of some/any matrix 𝑇 representing 𝜏 with respect to some
basis.
TRAIL: Let (𝑉, 𝐸) be a graph. A trail 𝑣0 𝑣1 𝑣2 . . . 𝑣𝑛 of length 𝑛 in the graph from a vertex
𝑎 to a vertex 𝑏 is a walk of length 𝑛 from 𝑎 to 𝑏 with the property that the edges 𝑣𝑖−1 𝑣𝑖
are distinct for 𝑖 = 1, 2, . . . , 𝑛. A trail in a graph is thus a walk in the graph which
traverses edges of the graph at most once.
TRANSCENDENTAL FUNCTION: A transcendental function is that function which is not
an algebraic function.
TRANSCENDENTAL NUMBER: A transcendental number is a number that cannot occur
as the root of a polynomial equation with rational coefficients. The transcendental
numbers are a subset of the irrational numbers. Most values for trigonometric functions
are transcendental, as is the number e. The number 𝜋 is transcendental. The square
roots of rational numbers are not transcendental, even though they are often irrational.
For example, 3 is a root of the equation 𝑥 2 − 3 = 0, so it is not transcendental.
TRANSFINITE NUMBER: A cardinal number relating to an infinite set such as aleph-null
and the cardinality of the continuum.
TRANSFORMATION GROUP: A group with elements that are transformations with the
binary operation of composition of transformations. For example, the set of rotations
about a fixed point form a transformation group, but the set of
reflections in axes through a fixed point do not, as the set is not closed under
composition.
TRANSITION MAPS: Charts in an atlas may overlap and a single point of a manifold may
be represented in several charts. If two charts overlap, parts of them represent the same
region of the manifold, just as a map of Europe and a map of Asia may both contain
Moscow. Given two overlapping charts, a transition function can be defined which goes
from an open ball in Rn to the manifold and then back to another (or perhaps the same)
open ball in Rn. The resultant map, like the map T in the circle example above, is called
a change of coordinates, a coordinate transformation, a transition function, or
a transition map.
TRANSITIVE PROPERTY: The transitive property of equality states that, if 𝑎 = 𝑏 and
𝑏 = 𝑐 , then 𝑎 = 𝑐 . All real and complex numbers obey this property. The transitive
property of inequality states that, if 𝑎 < 𝑏 and 𝑏 < 𝑐, then 𝑎 < 𝑐 . Real numbers obey
this property, but complex numbers do not.
TRANSITIVE RELATION: A binary relation ∼ on a set 𝑆 is transitive if, for 𝑎𝑙𝑙 𝑎, 𝑏 and 𝑐 in
𝑆, whenever 𝑎 ∼ 𝑏 and 𝑏 ∼ 𝑐 then 𝑎 ∼ 𝑐.
TRANSLATION (COMPLEX ANALYSIS): The map 𝑤 = 𝑧 + 𝛽 corresponds to a
translation. For, by this transformation, the figure in 𝑤 − 𝑝𝑙𝑎𝑛𝑒 is the same as figure in
𝑧 − 𝑝𝑙𝑎𝑛𝑒 with a different origin.

TRANSLATION: A translation occurs when we shift the axes of a Cartesian coordinate

system. (Keeping the orientation of the axes the same; otherwise there would be a
rotation.) If the new coordinates are called 𝑥 ′ and 𝑦 ′ , and the amount that the 𝑥 −axis is
shifted is 𝑕 and the amount that the 𝑦 −axis is shifted is 𝑘, then there is a simple relation
between the new coordinates and the old coordinates:
𝑥′ = 𝑥 + 𝑕
𝑦′ = 𝑦 + 𝑘

TRANSPORTATION PROBLEM: From the name itself, it is clear that transportation

problem means a problem where something is to be transferred.
Let us suppose we have a product which is to be transported from a number of centres
called ‘origin’ or ‘sources’ to a number of places called ‘destinations’. The costs of
transportation along different routes are different and are known.

The main objective is to minimize the cost of associated with such transportation from
the place of supply to their destination. These special types of linear programming
problems are called ‘Transportation problem’.

TRANSPOSE: The transpose of a matrix is formed by changing all the columns in the
original matrix into corresponding rows in the transposed matrix.
TRANSPOSED CONJUGATE OF A MATRIX: The transpose of the conjugate of a matrix A is
called transposed conjugate of A and 𝐵 is denoted by 𝐴° or A∗ .

Obviously the conjugate of the transpose of A is the same as the transpose of the
conjugate of A i.e., A′ = A ′ = 𝐴°.

If A = aij then A° = bji where bji = aij i.e., the (j, t)th element of A°= the
m×n, n×m

conjugate complex of the (i, j)th element of A.

For example, If

1 + 2i 2 − 3i 3 + 4i 1 + 2i 4 − 5i 8
′
A = 4 − 5i 5 + 6i 6 − 7i , then A = 2 − 3i 5 + 6i 7 + 8i
8 7 + 8i 7 3 + 4i 6 − 7i 7

1 − 2i 4 + 5i 8
and 𝐴′ = 𝐴° 2 + 3i 5 − 6i 7 − 8i .
3 − 4i 6 + 7i 7

If A° and B° be the transposed conjugates of A and 𝐵 respectively, then

(i) A° = A;
(ii) A + B ° = A° + B°, A and B being of the same size.
(iii) kA ° = k A°, 𝑘 being any complex number;
(iv) AB = B°A°, A and B being conformable to multiplication.

TRANSPOSE OF MATRIX: Let A = aij . Then the 𝑛 × 𝑚 matrix obtained from 𝐴 by

m×n

changing its rows into columns and its columns into rows is called the transpose of 𝐴
and is denoted by the symbol 𝐴′ or 𝐴𝑇 .
The operation of interchanging rows with columns is called transposition. Symbolically
if A = aij ,then A′ = bji , where bji = aij , i.e. the (𝑗, 𝑖)𝑡𝑕 element of A′ is the
m×n n×m

(𝑖. 𝑗)𝑡𝑕 element of A.

1 2 3 4
For example, the transpose of the 3 × 4 matrix A = 2 3 4 1 is the 4 × 3 matrix
3 4 2 1 3×4

1 2 3
2 3 4
A′ = .
3 4 2
4 1 1 4×3

The first row of A is the first column of A′. The second row of A is the second column of
A′ . the third row of A is the third column of A.

If A′ and 𝐵′ be the transpose of A and B respectively,then

(i) A′ = A
(ii) A + B ′ = A′ + B′ , A and B being of the same size.
(iii) kA ′ = kA′ , k being any complex number.
(iv) AB ′ = B′ A′ , A and B being conformable to multiplication.

TRANSPOSITION: Given a set 𝑋 = {𝑎1 , 𝑎2 , . . . , 𝑎𝑛 }, a transposition is a permutation

(bijective function of 𝑋 onto itself) 𝑓 such that there exist indices 𝑖, 𝑗 such that
𝑓(𝑎𝑖 ) = 𝑎𝑗 , 𝑓(𝑎𝑗 ) = 𝑎𝑖 and 𝑓(𝑎𝑘 ) = 𝑎𝑘 for all other indices 𝑘.
Example: If 𝑋 = {𝑎, 𝑏, 𝑐, 𝑑, 𝑒} the function 𝛼 given by
𝛼(𝑎) = 𝑎
𝛼(𝑏) = 𝑒
𝛼(𝑐) = 𝑐
𝛼(𝑑) = 𝑑
𝛼(𝑒) = 𝑏
is a transposition.
One of the main results on symmetric groups states that any permutation can be
expressed as composition of transpositions, and for any two decompositions of a given
permutation, the number of transpositions is always even or always odd.
TRANSVERSABLE GRAPH: A graph that can be drawn without removing pen from paper
or going over the same edge twice.
TRAPEZIUM: A quadrilateral with two parallel sides. If the parallel sides have lengths 𝑎
1
and 𝑏, and the distance between them is 𝑕, the area of the trapezium equals 2 𝑕 𝑎 + 𝑏 .

TRAPEZOIDAL FUZZY NUMBER: A fuzzy set 𝐴 is called trapezoidal fuzzy number with
tolerance interval [𝑎, 𝑏], left width 𝛼 and right width 𝛽 if its membership function has
the following form
𝑎−𝑡
1−;𝑎 − 𝛼 ≤ 𝑡 ≤ 𝑎
𝛼
1; 𝑎 ≤ 𝑡 ≤ 𝑏
𝐴 𝑡 = 𝑡−𝑏
1− ;𝑏 ≤ 𝑡 ≤ 𝑎 + 𝛽
𝛽
0; 𝑜𝑡𝑕𝑒𝑟𝑤𝑖𝑠𝑒

and we use the notation 𝐴 = (𝑎, 𝑏, 𝛼, 𝛽).

TRAVELLING SALESMAN PROBLEM: We consider a special type of sequencing problem

the travelling salesman problem. A travelling salesman wants to minimize the total
distance travelled (or time or money) during his visit of n cities. A similar problem arise
when 𝑛 items say 𝐴𝑖 , i= 1,2 …,n, are to be produced on a machine in continuation , given
that 𝑐𝑖𝑗 𝑖, 𝑗 = 1,2 … , 𝑛 is he set up cost of the machine when items 𝐴𝑖 is followed by 𝐴𝑗 .
Note that 𝑐𝑖𝑗 = ∞ 𝑤𝑕𝑒𝑛 𝑖 = 𝑗 i.e., we don’t produce the item 𝐴𝑖 again after 𝐴𝑗 . The
individual set up costs is arranged in the form of the adjacent square matrix.

∞ 𝑐12… 𝑐11 … 𝑐𝑖𝑛

𝑐21 ∞… 𝑐21… 𝑐2𝑛

⋮ ⋮ ⋮ ⋮

𝑐𝑖2 𝑐𝑖2 ∞ 𝑐𝑖𝑛

⋮ ⋮ ⋮ ⋮

𝑐𝑛1 𝑐𝑛2 … 𝑐𝑛𝑖 … ∞

Our problem is to determine a set of 𝑛 elements of this matrix, one of each row and one
in each column, so as to minimize the sum of the elements determined above. Here two
extra restrictions are imposed. One restriction is that we cannot select the element in
the leading diagonal as we have already assumed the elements of leading diagonal as we
have already assumed the elements of leading diagonal to be infinity. The other
restriction is that we don’t produce an item again until all the items are produced once.

A travelling salesman problem is said to be symmetric or asymmetric according as cost

matrix is symmetric or not. This it is symmetric if the cost from 𝐴𝑖 𝑡𝑜 𝐴𝑗 is the same as
that from 𝐴𝑖 𝑡𝑜 𝐴𝑗

A travelling salesman problem may be stated as follows:

Like assignment problem, if we represent “going” and “not going” of the salesman from
𝐴𝑖 𝑡𝑜 𝐴𝑗 station by saying 𝑥𝑖𝑗 = 1 𝑎𝑛𝑑 𝑥𝑖𝑗 = 0 respectively, then we wish to determine
𝑥𝑖𝑗 , 𝑖, 𝑗 = 1 … , 𝑛, which minimizes

𝑛 𝑛 𝑛
𝑧= 𝑖=1 𝑗 =1 𝑐𝑖𝑗 𝑥𝑖𝑗 s. t. 𝑗 =1 𝑥𝑖𝑗 = 1, 𝑖 = 1,2, … , 𝑛

𝑛
𝑗 =1 𝑥𝑖𝑗 = 1, 𝑗 = 1,2, … , 𝑛

𝑥𝑖𝑗 = 0 𝑜𝑟 1,

and one extra restriction that the 𝑥𝑖𝑗 must be so chosen hat no city is visited twice until
the tour of all the cities is completed. Note that as we have written ∞ in the leading
diagonal 𝑥𝑖𝑗 cannot be 1 when 𝑖 = 𝑗, because then 𝑧 will not be minimum. Thus the
second restriction of no going to 𝐴𝑖 again just after 𝐴𝑖 is automatically satisfied.

TREE: A connected graph with no cycles. It can be shown that a connected simple graph
with 𝑛 vertices is a tree if and only if it has 𝑛 – 1 edges. Note that a non-trivial tree
contains at least one pendant vertex.
TREE DIAGRAM: A tree diagram illustrates all of the possible results for a process with
several stages. Following figure illustrates a tree diagram that shows all of the possible
results for tossing three coins.
TREE TRAVERSALS: A tree traversal is an algorithm for visiting all the nodes in a rooted
tree exactly once. The constraint is on rooted trees, because the root is taken to be the
starting point of the traversal. A traversal is also defined on a forest in the sense that
each tree in the forest can be iteratively traversed (provided one knows the roots of
every tree beforehand). This entry presents a few common and simple tree traversals.
TRIANGLE: A triangle is a three-sided polygon. The three points where the sides
intersect are called vertices. Triangles are sometimes identified by listing their vertices,
as in triangle ABC.

TRIANGLE INEQUALITY (COMPLEX NUMBERS): If 𝑧1 and 𝑧2 are complex numbers, then

|𝑧1 + 𝑧2 | ≤ |𝑧1 | + |𝑧2 |. This result is known as the triangle inequality because it
follows from the fact that |𝑂𝑄| ≤ |𝑂𝑃| + |𝑃𝑄|, where 𝑂𝑃, 𝑃𝑄 and 𝑂𝑄 represent 𝑧1 , 𝑧2
and 𝑧1 + 𝑧2 in the complex plane.
TRIANGLE INEQUALITY (PLANE): For points 𝐴, 𝐵 and 𝐶 in the plane, |𝐴𝐶| ≤ |𝐴𝐵| +
|𝐵𝐶|. This result, the triangle inequality, says that the length of one side of a triangle is
less than or equal to the sum of the lengths of the other two sides.
TRIANGLE INEQUALITY (VECTORS): Let |𝒂| denote the length of the vector 𝒂. For
vectors 𝒂 and 𝒃, |𝒂 + 𝒃| ≤ |𝒂| + |𝒃|. This result is known as the triangle inequality,
since it is equivalent to saying that the length of one side of a triangle is less than or
equal to the sum of the lengths of the other two sides.
TRIANGLE OF FORCES (MECHANICS): If three forces act at a point on a body in
equilibrium their vector sum must be zero, and consequently the forces can be drawn in
a closed triangle.
TRIANGULAR FUZZY NUMBER: A fuzzy set 𝐴 is called triangular fuzzy number with peak
(or center) 𝑎, left width 𝛼 > 0 and right width 𝛽 > 0 if its membership function has
the following form
𝑎−𝑡
1− ;𝑎 − 𝛼 ≤ 𝑡 ≤ 𝑎
𝛼
𝐴 𝑡 = 𝑡−𝑎
1− ;𝑎 ≤ 𝑡 ≤ 𝑎 + 𝛽
𝛽
0; 𝑜𝑡𝑕𝑒𝑟𝑤𝑖𝑠𝑒

and we use the notation 𝐴 = (𝑎, 𝛼, 𝛽).

TRIANGULAR INEQUALITY: It states that for any triangle, the sum of the lengths of any
two sides must be greater than or equal to the length of the remaining side. If 𝑥, 𝑦
and 𝑧 are the lengths of the sides of the triangle, then the triangle inequality states that

with equality only in the degenerate case of a triangle with zero area. In Euclidean
geometry, the triangle inequality is a theorem about distances, and it is written using
vectors and vector lengths:

𝑥+𝑦 ≤ 𝑥 + 𝑦 .
TRIANGULARIZATION THEOREM OR JACOBI’S THEOREM: Every square matrix is
unitarily similar to a triangular matrix.
TRIANGULATION: The method of fixing a point by using directions from two known
points to construct a triangle.
TRIANGULAR MATRIX: A square matrix that is either lower triangular or upper
triangular. It is lower triangular if all the entries above the main diagonal are zero, and
upper triangular if all the entries below the main diagonal are zero.
𝑛(𝑛+1)
TRIANGULAR NUMBER: An integer of the form , where 𝑛 is a positive integer. The
2

first few triangular numbers are 1, 3, 6, 10, 15, … …..

TRIDIAGONAL MATRIX: A square matrix with zero entries everywhere except on the
main diagonal and the neighbouring diagonal on either side.
TRIGONOMETRY: The area of mathematics relating to the study of trigonometric
functions in relation to the measurements in triangles.
TRILLION: A million million (1012 ).
TRINOMIAL: A trinomial is the indicated sum of three monomials. For example,
8 + 23𝑥 2 − 43𝑎4 𝑏 3 is a trinomial.
TRIPLE INTEGRAL: A triple integral means to integrate a function over an entire volume.
For example, if 𝑟(𝑥, 𝑦, 𝑧) represents the density of matter at a point (𝑥, 𝑦, 𝑧), then
𝑧=𝑐 𝑦 =𝑏 𝑥=𝑎

𝑟 𝑥, 𝑦, 𝑧 𝑑𝑥𝑑𝑦𝑑𝑧
𝑧=0 𝑦 =0 𝑥=0

gives the total mass contained in the parallelepiped from 𝑥 = 0 to 𝑥 = 𝑎, 𝑦 = 0 to

𝑦 = 𝑏, and 𝑧 = 0 to 𝑧 = 𝑐.
TRISECT: To trisect an object means to cut it in three equal parts.
TRIVIAL GRAPH: A graph is said to be trivial if it consists of a single vertex.
TRIVIAL PATH: A path in a graph is said to be trivial if it is a path v of length zero
determined by a single vertex v of v; otherwise it is said to be non-trivial.
TRIVIAL SOLUTION: The solution of a homogeneous set of linear equations in which all
the unknowns are equal to zero.
TRIVIAL SUBGROUP: The subset of a group which contains only the identity element. It
will always be a group itself.
TRIVIAL TRAIL: A trail in a graph is said to be trivial if it is a trail v of length zero
determined by a single vertex v of v; otherwise it is said to be non-trivial.
TRIVIAL WALK: A walk in a graph is said to be trivial if it is a walk v of length zero
determined by a single vertex v of v; otherwise it is said to be non-trivial.
TRUNCATED CONE: A truncated cone consists of the section of a cone between the base
and another plane that intersects the cone between the base and the vertex. It looks like
a cone whose top has been chopped off.
TRUNCATED PYRAMID: A truncated pyramid consists of the section of a pyramid
between the base and another plane that intersects the pyramid between the base and
the vertex. It looks like a pyramid whose top has been chopped off.
TRUNCATION: The truncation of a number is found by dropping the fractional part of
that number. It is equal to the largest integer that is less than or equal to the original
number. For example, the truncation of 65.94 is equal to 65.
TRUNCATION ERROR: The error resulting in using an approximation obtained by
truncation. Truncation error is the difference between a truncated value and the actual
value. A truncated quantity is represented by a numeral with a fixed number of allowed
digits, with any excess digits "chopped off" (hence the expression "truncated").

As an example of truncation error, consider the speed of light in a vacuum. The official
value is 299,792,458 meters per second. In scientific (power-of-10) notation, that
quantity is expressed as 2.99792458 x 108. Truncating it to two decimal places yields
2.99 x 108. The truncation error is the difference between the actual value and the
truncated value, or 0.00792458 x 108. Expressed properly in scientific notation, it is
7.92458 x 105.

In computing applications, truncation error is the discrepancy that arises from

executing a finite number of steps to approximate an infinite process. For example, the
infinite series 1/2 + 1/4 + 1/8 + 1/16 + 1/32 ... adds up to exactly 1. However, if we
truncate the series to only the first four terms, we get 1/2 + 1/4 + 1/8 + 1/16 = 15/16,
producing a truncation error of 1 - 15/16, or 1/16.

TRUTH TABLE: A truth table is a table showing whether a compound logic sentence will
be true or false, based on whether the simple sentences contained in the compound
sentence are true. Each row of the table corresponds to one set of possible truth values
for the simple sentences.
TRUTH VALUE: In logic, a sentence is assigned one of two truth values. One of the truth
values is labeled T, or 1; it corresponds to “true.” The other truth value is labeled F, or 0;
it corresponds to “false.”
t-TEST: A test to determine whether or not a sample of size 𝑛 with mean comes from a
normal distribution with mean 𝜇. The statistic 𝑡 given by
𝑛 𝑥−𝜇
𝑡=
𝑠
has a 𝑡-distribution with 𝑛 – 1 degrees of freedom, where 𝑠 2 is the (unbiased) sample
variance. The 𝑡-test can also be used to test whether a sample mean differs significantly
from a population mean, and to test whether two samples have been drawn from the
same population.
TURNING POINT: A point on the graph 𝑦 = 𝑓(𝑥) at which 𝑓′(𝑥) = 0 and 𝑓′(𝑥) changes
sign. A turning point is either a local maximum or a local minimum.
TWIN PRIMES: A pair of prime numbers that differ by 2. For example, 11 and 13, 29 and
31, 71 and 73 are twin primes. A conjecture that there are infinitely
many such pairs has been neither proved nor disproved.

TWO DIMENSIONAL DOUBLET: A combination of a source and a sink of equal strength

𝑚 at a distance 𝛿𝑠 apart where 𝛿𝑠 is taken to be infinitely small and 𝑚 infinitely great,
such that the product 𝑚 𝛿𝑠 is finite (let it be equal to 𝑢) is called a two dimensional
doublet of strength 𝑢. The line 𝛿𝑠 drawn from sink to source is called the axis of the
doublet.

The function (1/𝑧) has a singularity at 𝑧 = 0. This singularity is called a doublet.

TWO-PERSON ZERO-SUM GAME: A game with two players in which the total payoff is
zero, i.e. anything which one player gains is directly at the expense of the other player.
TWO-PHASE SIMPLEX METHOD:
Initialization Modify constraints of innovative problem to achieve an obvious Basic
feasible solution for the artificial problem.
Phase I Discover a Basic feasible solution for the actual problem by minimizing
the sum of the artificial variables.
Phase II Leave the artificial variables, replace with phase 2 objective function.
Re-establish appropriate structure. Locate an optimal solution for the
real problem.
TWO-SAMPLE TESTS (STATISTICS): Any test which is to be applied to two independent
samples, in contrast to paired-sample tests when the two samples are combined before
applying a one-sample test to the resulting sample.
TWO-TAILED TEST: In a two-tailed test the critical region consists of both tails of a
distribution. The null hypothesis is rejected if the test-statistic value is either too large
or too small.
TYPE m FUZZY SET: A type m fuzzy set is a fuzzy set whose membership values are type
𝑚 − 1, 𝑚 > 1, fuzzy sets on [0, 1].
TYPES OF MODELS:

ICONIC MODELS: Iconic models represent the system as it is but in different size. Thus,
ionic models are obtained by enlarging or reducing the size of the system. In other
words, they are images.

Some common examples are photographs, drawing, model aeroplanes, ships, engines,
globes, maps etc. A toy aeroplane is an ionic model of a real one.

ANALOGUE MODEL: It represents a system or an object of an inquiry by utilizing a set of

properties different from what the original system process. For example graphs are
very simple analogues because distance is used to represent the properties such as
time, number, percent, age weigh and many other properties. Contour lines on a map
represent the rise and fall of the heights. In general, analogues are less specific, less
concrete but easier to manipulate than are iconic models.

MATHEMATICAL (SYMBOLIC) MODEL: The symbolic or mathematical model is one

which employs a set of mathematical symbols (i.e., letters, numbers, etc) to represent
decision variables of the system. These variables are related together by means of a
mathematical equation or a set of equations to describe the behavior of the system. The
solution of the problem is then obtained by applying well-developed mathematical
techniques to the model.

The symbolic model is usually the easiest to manipulate experimentally and it is most
general and abstract. Its function is more often explanatory rather than descriptive.

DESCRIPTIVE MODELS: A descriptive model simply describes some aspects of a

situation based on observation, survey, questionnaire results, or other available data.
The result of an opinion poll represents a descriptive model.

PREDICTIVE MODEL: Such models can answer ‘what if’ type of question, i.e, they can
make predictions regarding certain event. For example, based on the survey results,
television networks such models attempt to explain and predict the election result
before all the votes are actually counted.

PRESCRIPTIVE OR NORMALATIVE MODELS: When a predictive model has been

repeatedly successful, it can be used to prescribe a source of action. Linear
programming is a normalated or prescriptive model because it prescribes what the
managers ought to do.

DETERMINISTIC MODELS: Such models assume conditions of conditions of complete

certainty and perfect knowledge. For example, linear programming, transportation and
assignment models are deterministic models.

PROBABILISTIC MODELS: These models handle those situations in which the

consequences or pay off of managerial actions can not be predicted with certainty.
However, it is possible to forecast a pattern of events, based on which managerial
decisions can be made. For example, insurance companies are willing to insure against
risk of fire, accidents, sickness and so on because the patterns of events have been
compiled in the form of probability distributions.

DYNAMIC MODELS: In these models, time is considered as one of the important variable
and admits the impact of changes generated by time. Also, in dynamic models not only
one, but a series of interdependent is required during the planning horizon.

STATIC MODELS: These models do not consider the impact of changes that take place
during the planning horizon, i.e. they are independent of time. Also, in static model only
one decision is needed for the duration of a given time period.

TYPE 1 ERROR: A type 1 error occurs when the null hypothesis is rejected when it is
actually true.
TYPE 2 ERROR: A type 2 error occurs when the null hypothesis is accepted when it is
actually false.
ULTRAFILTER: A filter 𝐹 on 𝑀 is called ultrafilter if for every 𝐴 ⊆ 𝑀 either
𝐴 ∇ 𝐹 𝑜𝑟 𝑀 − 𝐴 ∇ 𝐹. Ultrafilters are precisely the maximal filters with respect to
inclusion. Using AC, it can be shown that every system of sets, which has finite
intersection property, is contained in an ultrafilter.
ULTRAFILTER LEMMA: Ultrafilter lemma states that every filter on a set 𝑋 is a subset of
some ultrafilter on 𝑋. This lemma is most often used in the study of topology. An
ultrafilter that does not contain finite sets is called non-principal. The ultrafilter lemma,
and in particular the existence of non-principal ultrafilters follows easily from Zorn's
lemma.

UNARY OPERATION: A unary operation on a set 𝑆 is a rule that associates with any
element of 𝑆 a resulting element. If this resulting element is always also in 𝑆, then it is
said that 𝑆 is closed under the operation. The following are examples of unary
operations: the rule that associates with each integer 𝑎 its negative – 𝑎; the rule that
associates with each non-zero real number 𝑎 its inverse 1/𝑎; and the rule that
associates with any subset 𝐴 of a universal set 𝐸 its complement 𝐴′.
UNBALANCED ASSIGNMENT PROBLEM: Sometimes, during a problem, the number of
jobs is not equal to number of persons, that means our assignment matrix is not square.
Such kinds of problems are known as unbalanced assignment problem.

There types of problems are solved by the introduction of a supposed or dummy rows
or columns in which each element has zero cost to form a square matrix. The new
square matrix can be solved by the usual method of balanced assignment problem.

UNBOUNDED DUAL THEOREM: Let X 0 be an open set in Rn and let 𝜃 and 𝑔 be

differentiable on X 0 . If there exists a dual feasible point x1 , u1 such that

g x1 ) + ∆g(x1 z ≤ 0

has no solution on z ∇ Rn , then the dual problem has an unbounded objective function
on the set by dual feasible points Y.

UNBOUNDED FUNCTION: A function 𝑓(𝑥) not possessing both an upper and a lower
bound. So for all positive real values 𝑀 there is a value of the independent variable x for
which |𝑓(𝑥)| > 𝑀. For example, 𝑓(𝑥) = 𝑥 4 is unbounded because 𝑓(𝑥) ≥ 0 but
𝑓(𝑥) → ∞ as 𝑥 → ±∞, i.e. it is bounded below but not above, while 𝑓(𝑥) = 𝑥 3 has
neither upper nor lower bound.
UNBOUNDED SOLUTION: The solution of a LPP is said to be unbounded solution or LPP
is said to possess an unbounded solution if the value of the objective function can be
increased or decreased indefinitely.
UNCONDITIONAL INEQUALITY: An inequality that is always true, irrespective of the
values taken by the variables so 𝑥 2 + 4𝑥 + 13 > 0 is unconditional since it can be
expressed as (𝑥 + 2)2 + 9 > 0 which is always > 0, but 𝑥 2 + 4𝑥 + 3 > 0 is a
conditional inequality because it can be expressed as (𝑥 + 2)2 > 1 which is only true
for 𝑥 > −1 or 𝑥 < – 3.
UNCOUNTABLE SET: A set which is not countable, so the elements cannot be put into
one to-one correspondence with the natural numbers, or a subset of them.
UNDEFINED TERM: An undefined term is a basic concept that is described, rather than
given a rigorous definition. It would be impossible to rigorously define every term,
because sooner or later the definitions would become circular. “Plane” is an example of
an undefined term from spherical geometry.
UNDIRECTED GRAPH: An undirected graph (𝑉, 𝐸) consists of a finite set 𝑉 together with
a subset 𝐸 of 𝑉2 (where 𝑉2 is the set consisting of all subsets of 𝑉 with exactly two
elements). The elements of 𝑉 are the vertices of the graph; the elements of 𝐸 are the
edges of the graph.

UNIFORM AND NON-UNIFORM FLOWS: If at every point the velocity vector is identical
in magnitude and direction at any given instant, or, the conditions and properties are
independent of the coordinate of the direction in which the fluid is moving then the
motion is said to be uniform. If the flow characteristics at any given time change with
distance, it is said to be non-uniform. For example, a liquid pumped through a long
straight pipe has uniform flow whereas a liquid flowing through a curved pipe has non-
uniform flow.

UNIFORM BOUNDEDNESS PRINCIPLE: Let 𝑋 be a Banach space and 𝑌 be a normed

vector space. Suppose that 𝐹 is a collection of continuous linear operators from 𝑋 to 𝑌. If
for all 𝑥 in 𝑋 one has

then

UNIFORM DISTRIBUTION: The uniform distribution on the interval [𝑎, 𝑏] is the

continuous probability distribution whose probability density function 𝑓 is given by
𝑓(𝑥) = 1/(𝑏 – 𝑎), where 𝑎 ≤ 𝑥 ≤ 𝑏. It has mean (𝑎 + 𝑏)/2 and variance
(𝑏 – 𝑎)2 /12. There is also a discrete form: on the set 1, 2, … , 𝑛, it is the probability
distribution whose probability mass function is given by 𝑃𝑟(𝑋 = 𝑟) = 1/2, for
𝑟 = 1, 2, … , 𝑛. For example, the random variable for the winning number in a lottery
has a uniform distribution on the set of all the
numbers entered in the lottery.
UNIFORM GRAVITATIONAL FORCE: A gravitational force, acting on a particular body
that is independent of the position of the body. The gravitational force on a body in a
limited region near the surface of a planet is approximately uniform.
UNIFORMLY CONTINUOUS FUNCTION: If the domain and range of a function are both
metric spaces in which, for every 𝜀 > 0 in the domain there exists a single 𝛿 > 0 in the
range so that for all 𝑥 in the domain
𝑓 𝑥 + 𝛿 –𝑓 𝑥 < 𝜀.
For continuity only, the value of 𝛿 can depend on x as well as 𝜀.
UNION: The union of two sets A and B (written as 𝐴 ∪ 𝐵) is the set of all elements that
are either members of A or members of B, or both. For example, the union of the sets
𝐴 = {2,3,4} and 𝐵 = {2,4,6,8} is the set 𝐴 ∪ 𝐵 = {2,3,4,6,8}. The union of the set of
whole numbers and the set of negative integers is the set of all integers.
UNION OF FUZZY SETS: The union of fuzzy sets A and B is a fuzzy set in U, denoted by A
U B, whose membership function is defined as
𝜇𝐴𝑈𝐵 (𝑥) = 𝑚𝑎𝑥[𝜇𝐴 (𝑥), 𝜇𝐵 (𝑥)]
UNIQUENESS OF ANALYTIC CONTINUATION: There cannot be more than one
continuation of analytic continuation 𝑓2 (𝑧) into the same domain.

UNIQUENESS OF MULTIPLICATIVE IDENTITIES: If (𝑀, . ) is an object with multiplication

and 1, 1′ ∇ 𝑀 are identities in the sense that 1𝑎 = 𝑎1 = 𝑎 and 1′ 𝑎 = 𝑎1′ = 𝑎 for all
𝑎 ∇ 𝑀, then 1 = 1′ .
UNIQUENESS THEOREM: A curve is uniquely determined, except as to position in space,
when its curvature and torsion are given functions of its arc length 𝑠.
UNITARILY SIMILAR MATRICES: Let 𝐴 and 𝐵 square matrices of order 𝑛. Then 𝐵 is said
to be unitarily similar to 𝐴 if there exists unitary matrix 𝑃 such that 𝐵 = 𝑃−1 𝐴𝑃.

 If 𝐴 and 𝐵 are unitarily similar, then they are similar also.

 The relation of being ‘unitarily similar’ is an equivalence relation in the set of all
𝑛 × 𝑛 matrices over the field of complex numbers.
  Every Hermitian matrix is unitarily similar to a diagonal matrix.
 An 𝑛 × 𝑛 Hermitian matrix 𝐻 has 𝑛 mutually orthogonal eigenvectors is the
complex vector space Vn .
 Any two eigenvectors corresponding to two distinct eigenvalues of a Hermitian
matrix are orthogonal.

UNITARY MATRIX: A square matrix 𝐴 is said to be unitary if 𝐴° 𝐴 = 𝐼. If 𝐴, 𝐵 be 𝑛- rowed

unitary matrices, AB and BA are also unitary matrices.

A unitary matrix over the field of real numbers in orthogonal i.e., a real unitary matrix is
an orthogonal matrix.

UNITARY OPERATOR: 𝑈: 𝐻 → 𝐻 is said to be a unitary operator on a Hilbert space 𝐻 if

𝑈 ∗ = 𝑈 −1 , i.e. 𝑈 ∗ 𝑈 = 𝑈𝑈 ∗ = 𝐼.

For an operator U on a complex Hilbert space 𝐻 the following are equivalent:

1. 𝑈 is unitary;

2. 𝑈 is surjection and an isometry, i.e. ||𝑈𝑥|| = ||𝑥|| 𝑓𝑜𝑟 𝑎𝑙𝑙 𝑥 ∇ 𝐻;

3. 𝑈 is onto and preserves the inner product, i.e. ⟨ 𝑈𝑥, 𝑈𝑦 ⟩ = ⟨ 𝑥, 𝑦 ⟩ for all 𝑥, 𝑦 ∇ 𝐻.

UNIT CIRCLE: In the plane, the unit circle is the circle of radius 1 with its centre at the
origin. In Cartesian coordinates, it has equation 𝑥 2 + 𝑦 2 = 1. In the complex plane, it
represents those complex numbers 𝑧 such that |𝑧| = 1.
UNIT ELEMENT: An element 𝑢 of an integral domain 𝑅 is said to be a unit if there exists
some element 𝑢−1 of 𝑅 such that 𝑢𝑢−1 = 1.

UNITS MATRIX OR IDENTITY MATRIX: A square matrix each of whose diagonal elements
is 1 and each of whose non-diagonal elements is equal to zero is called a unit matrix or
an identity matrix and is denoted by I. In will denote a unit matrix of order n. Thus a
square matrix 𝐴 = 𝑎𝑖𝑗 is a unit matrix if𝑎𝑖𝑗 = 1 when 𝑖 = 𝑗 and 𝑎𝑖𝑗 = 0 when 𝑖 ≠ 𝑗.

For example,

1 0 0 0
0 1 0 0
I4 =
0 0 1 0
0 0 0 1
 1 0 0
1 0
I3 0 1 0 , I2
0 1
0 0 1

Are unit matrices of orders 4,3,2 respectively.

UNIT VECTOR: A vector, whose modulus is unity, is called a unit vector. The unit vector
in the direction of vector 𝜶 is represented by 𝜶. It is read as ‘𝜶 Cap’.

UNIVALENT FUNCTION: The map ω = f(z) is called one-one (or one to one) function if

f z1 ⇒ z1 = z2

or z1 ≠ z2 ⇒ f z1 ≠ f z2

If the function ω = f(z) one-one, then it is also called univalent function, a fancy term for
one-one map.

A function f(z) is called univalent at 𝑧 = 𝑧0 if it is univalent in a neighbourhood of 𝑧0 .

UNIVERSAL QUANTIFIER: An upside-down letter A, ∀is used to represent the

expression “For all . . . ,” and is called the universal quantifier. For example, if 𝑥 is
allowed to take on real-number values, then the sentence “For all real numbers, the
square of the number is nonnegative” can be written as
∀𝑥 , 𝑥 2 ≥ 0
UNIVERSAL SET: A universal set is the collection of all objects in a particular context or
theory. All other sets in that framework constitute subsets of the universal set, which is
denoted as an uppercase italic letter U. The objects themselves are known as elements
or members of U. The precise definition of U depends the context or theory under
consideration. For example, we might define U as the set of all living things on planet
earth. In that case, the set of all dogs is a subset of U, the set of all fish is another subset
of U, and the set of all trees is yet another subset of U. If we define U as the set of all
animals on planet earth, then the set of all dogs is a subset of U, the set of all fish is
another subset of U, but the set of all trees is not a subset of U. Some philosophers have
attempted to define U as the set of all objects (including all sets, because sets are
objects). This notion of U leads to a contradiction, because U, which contains
everything, must therefore contain the set of all sets that are not members of
themselves. In 1901, the philosopher and logician Bertrand Russell proved that this
state of affairs leads to a paradox. Today, mathematicians and philosophers refer to this
result as Russell's Paradox.

UNIT CIRCLE: A unit circle is a circle with radius 1. If the unit circle is centered at the
origin, and (𝑥, 𝑦) is a point on the circle such that the line from the origin to that point
makes an angle 𝑢 with the 𝑥 −axis, then 𝑠𝑖𝑛 𝑢 = 𝑦 and 𝑐𝑜𝑠 𝑢 = 𝑥.
UNIT VECTOR: A unit vector is a vector of length 1. It is common to use 𝒊 to represent
the unit vector along the 𝑥-axis—that is, the vector whose components are (1,0,0).
Likewise, 𝒋 is used to represent (0,1,0), and 𝒌 represents (0,0,1). A three-dimensional
vector whose components are (𝑥, 𝑦, 𝑧) can be written as the vector sum of each
component times the corresponding unit vector: 𝑥, 𝑦, 𝑧 = 𝑥𝒊 + 𝑦𝒋 + 𝑧𝒌.
UNKNOWN CONSTANT: A constant whose value is not currently known. For example, in
the general equation of a straight line 𝑦 = 𝑚𝑥 + 𝑐, 𝑥 and 𝑦 are variables and 𝑚 and 𝑐
are unknown constants which are determined by the gradient and intercept of the
particular line.
UPPER AND LOWER BOUNDS: Characteristics of sets on the real line. The least upper
bound of a given set of real numbers is the smallest number bounding this set from
above; its greatest lower bound is the largest number bounding it from below. This will
now be restated in more detail. Let there be given a subset 𝑿 of the real numbers. A
number 𝜷 is said to be its least upper bound, denoted by sup𝑿 (from the Latin
"supremum" — largest), if every number 𝒙 ∇ 𝑿 satisfies the inequality 𝒙 ≤ 𝜷, and if for
any 𝜷′ < 𝜷 there exists an 𝒙′ ∇ 𝑿 such that 𝒙′ > 𝜷′. A number 𝜶 is said to be the greatest
lower bound of 𝑿, denoted by 𝒊𝒏𝒇𝑿 (from the Latin "infimum"-smallest), if every 𝒙 ∇
𝑿 satisfies the inequality 𝒙 ≥ 𝜶, and if for any 𝜶′ > 𝜶 there exists an 𝒙′ ∇ 𝑿 such
that 𝒙′ < 𝜶′.

Examples.

𝑖𝑛𝑓[𝑎, 𝑏] = 𝑎, 𝑠𝑢𝑝[𝑎, 𝑏] = 𝑏;

𝑖𝑛𝑓(𝑎, 𝑏) = 𝑎, 𝑠𝑢𝑝(𝑎, 𝑏) = 𝑏;

if the set 𝑋 consists of two points 𝑎 and 𝑏, 𝑎 < 𝑏, then

𝑖𝑛𝑓𝑋 = 𝑎, 𝑠𝑢𝑝𝑋 = 𝑏.
UPPER BOUND FOR A COMPLEX INTEGRAL: If a function 𝑓(𝑧) is continuous on a
contour 𝐶 of length 𝑙 𝑎𝑛𝑑 𝑖𝑓 𝑀 be the upper bound of 𝑓(𝑧) on 𝐶, then ∫𝐶 𝑓(𝑧)𝑑𝑧 ≤
𝑀𝑙.

UPPER RIEMANN INTEGRAL: Let 𝑓 be a bounded function on a bounded interval [𝑎, 𝑏].
Then upper Riemann integral of 𝑓 over [𝑎, 𝑏] is the infimum of 𝑈 𝑃, 𝑓 over all
𝑏
partitions 𝑃 of [𝑎, 𝑏] and is denoted as ∫𝑎 𝑓𝑑𝑥.
UPPER RIEMANN-STIELTJES SUM: Let 𝑓 be a bounded real valued function on a
bounded interval [𝑎, 𝑏] and let 𝑔 be a monotonically non-decreasing function on [𝑎, 𝑏].
Let 𝑃 = 𝑥0 , 𝑥1 , 𝑥2 , 𝑥3 , … , 𝑥𝑛 be a partition of [𝑎, 𝑏]. We write
∆𝑔𝑖 = 𝑔𝑖 − 𝑔𝑖−1
Then upper Riemann-Stieltjes sum denoted by 𝑈(𝑃, 𝑓, 𝑔) is defined as
𝑛

𝑈 𝑃, 𝑓, 𝑔 = 𝑀𝑖 ∆𝑔𝑖
𝑖=1

Where
𝑀𝑖 = sup 𝑓(𝑥)
𝑥∇[𝑥 𝑖−1 ,𝑥 𝑖 ]

UPPER RIEMANN SUM: Let 𝑓 be a bounded function on a bounded interval [𝑎, 𝑏]. Let
𝑃 = 𝑥0 , 𝑥1 , 𝑥2 , 𝑥3 , … , 𝑥𝑛 be a partition of [𝑎, 𝑏]. Then upper Riemann sum denoted by
𝑈(𝑃, 𝑓) is defined as
𝑛

𝑈 𝑃, 𝑓 = 𝑀𝑖 𝑥𝑖 − 𝑥𝑖−1
𝑖=1

Where
𝑀𝑖 = sup 𝑓(𝑥)
𝑥∇[𝑥 𝑖−1 ,𝑥 𝑖 ]

UPPER TRIANGULAR MATRIX: A square matrix A = aij is called an upper triangular

matrix if aij = 0 whenever i > 𝑗.

Thus in an upper triangular matrix all the elements below the principle diagonal are
zero.
 a11 a12 a13 … a1n
0 a22 a23 … a2n
For example A = 0 0 a33 … a3n
… … … … …
0 0 0 … amn

is an upper triangular matrix of the type n × n. Similarly

1 2 4 2
2 −9 0
0 3 −1 0
A= ,B = 0 1 2
0 0 2 1
0 0 1 3×3
0 0 0 8 4×4

is the upper triangular matrices.

URYSOHN’S LEMMA: Let 𝐾 be a compact space, and let 𝐸, 𝐹 be the closed subsets of 𝐾
with 𝐸 ∩ 𝐹 = ∅. There exists 𝑓: 𝐾 → [0,1] continuous with 𝑓(𝑥) = 1 for 𝑥 ∇ 𝐸 and
𝑓(𝑥) = 0 for 𝑥 ∇ 𝐹 (written 𝑓(𝐸) = {1} and 𝑓(𝐹) = {0}).

URYSOHN'S METRIZATION THEOREM: A topological space is separable and metrizable if

and only if it is regular, Hausdorff and second-countable.
VALUATION RING: Let 𝐾 be a field, and let 𝜈: 𝐾 → 𝑍 ∪ {∞} be a discrete valuation on 𝐾.
The valuation ring determined by this valuation is the subring 𝑅 of 𝐾 defined such that
𝑅 = {𝑐 ∇ 𝐾 ∶ 𝜈(𝑐) ≥ 0}. Let 𝐾 be a field, and let 𝜈: 𝐾 → 𝑍 ∪ {∞} be a discrete
valuation on 𝐾. Then the valuation ring 𝑅 of 𝜈 is integrally closed.
VAN DER CORPUT CONJECTURE: Suppose that 𝑠𝑖 ; 𝑖 ∇ 𝑁 is a real sequence in
𝐼 = [0, 1). Corresponding to any arbitrarily large real number 𝜅, there exist a positive
integer 𝑛 and two subintervals 𝐼1 and 𝐼2 , of equal length, of 𝐼 such that |𝑍(𝐼1 , 𝑛) −
𝑍(𝐼2 , 𝑛)| > 𝜅. In short, this conjecture expresses the fact that no sequence can, in a
certain sense, be too evenly distributed. This conjecture is true, as shown by van
Aardenne-Ehrenfest in 1945.
VAN DER WAERDEN THEOREM: Given any natural numbers 𝑡 and 𝑟, there exists 𝑁0 (𝑡, 𝑟)
such that for every natural number 𝑛 > 𝑁0 (𝑡, 𝑟), every partition of the set {1, 2, . . . , 𝑛}
into r subsets will yield a subset which contains 𝑡 terms in arithmetic progression.
VARIABLE: A variable is a symbol that is used to represent a value from a particular set.
For example, in algebra it is common to use letters to represent values from the set of
real numbers.
VARIANCE: The variance of a random variable X is defined to be
 𝑉𝑎𝑟 𝑋 = 𝐸[(𝑋 − 𝐸(𝑋)) × (𝑋 − 𝐸(𝑋))]
= 𝐸[(𝑋 − 𝐸(𝑋))2 ]
where E stands for “expectation.” The variance can also be found from the formula:
𝑉𝑎𝑟 𝑋 = 𝐸 𝑋 2 − [𝐸(𝑋)]2
The variance is often written as 𝜍 2 . (The Greek lowercase letter sigma (𝜍), is used to
represent the square root of the variance, known as the standard deviation.) The
variance is a measure of how widespread the observations of X are likely to be. If you
know for sure what the value of X will be, then 𝑉𝑎𝑟 𝑋 = 0.
VARIATION OF A CHARGE: The variation of a charge on a set 𝐴
𝑖𝑠 | 𝜈 |(𝐴) = 𝑠𝑢𝑝 𝑘 | 𝜈(𝐴𝑘 ) | for all disjoint splitting 𝐴 = ⋃𝑘 𝐴𝑘 .

1. If 𝜈 = µ1 − µ2 , then | 𝜈 |(𝐴) ≤ µ1 (𝐴) + µ2 (𝐴). The inequality becomes an identity

for disjunctive measures on 𝐴 (that is there is a partition of A i.e. 𝐴 = 𝐴1 ⊔ 𝐴2
such that µ2 (𝐴1 ) = µ1 (𝐴2 ) = 0).
2. For any charge 𝜈 the function | 𝜈 | is a ς-additive measure.
3. For any charge 𝜈 the function | 𝜈 | − 𝜈 is a 𝜍 −additive measure as well.
4. The collection of all charges on a 𝜍 −algebra 𝑅 is a linear space which is
𝑠𝑢𝑝
complete with respect to the distance: 𝑑 𝜈1 , 𝜈2 = 𝜈 𝐴 − 𝜈2 𝐴
𝐴∇𝑅 1

VECTOR AS A LINEAR COMBINATION OF VECTORS: A vector X which can be expressed

in the form: X = k1 X1 + k 2 X2 + ⋯ + k r Xr , is said to be a linear combination of the set of
vectors X1 , X2 , … , Xr . Here k1 , k 2 , … , k r are any numbers.

The following two results are quite obvious:

(i) If a set of vector is linearly dependent, then at least one member of the set can be
expressed as a linear combination of the remaining members.
(ii) If a set of vectors is linearly independent then no member of the set can be
expressed as a linear combination of the remaining member.

VECTOR PRODUCT OF FOUR VEECTORS: Let 𝒂, 𝒃 𝒄, 𝒅 be four vectors. Consider the

vector product of the vectors 𝒂 × 𝒃 and 𝒄 × 𝒅. This product can be written as
𝒂 × 𝒃 × 𝒄 × 𝒅 and is called the vector product of four vectors. It is a vector
perpendicular to 𝒂 × 𝒃 and, therefore coplanar with 𝒂 and 𝒃. Similarly it is a vector
coplanar with 𝒄 and 𝒅. Hence this vector must be parallel to the line of intersection of a
plane parallel to 𝒂 and 𝒃 with another plane parallel to 𝒄 and 𝒅.
VECTOR OF A FORCE: A vector, whose modulus is the magnitude of the given force and
whose direction is the same as that of the given force , is called the vector of the given
force or Force vector. The vector of a force will be denoted by the same symbol as
denotes the force. A force is completely specified by its vector and its point of
application.

VECTOR SPACES: A vector space over a field K is a set V which has two basic operations,
addition and scalar multiplications, satisfying certain requirements. Thus for every pair
𝑢, 𝑣 𝜖 𝑉, 𝑢 + 𝑣 𝜖 𝑉 is defined, and for every 𝛼 𝜖 𝐾 , 𝛼𝑣 𝜖 𝑉 is defined.For V to be called
vector space, the following axioms must be satisfied for all α, β ϵ K and all 𝑢, 𝑣 𝜖 𝑉.

(i) Vector addition satisfies axioms of a group.

(ii) 𝛼 𝑢 + 𝑣 = 𝛼𝑢 + 𝛼𝑣
(iii) 𝛼 + 𝛽 𝑣 = 𝛼𝑣 + 𝛽𝑣
(iv) 𝛼𝛽 𝑣 = 𝛼 𝛽𝑣
(v) 1. 𝑣 = 𝑣
Elements of the field 𝐾 will be called scalars and the elements of vector space V
are called vectors.

Examples of Vector Spaces:

1: 𝐾 𝑛 = 𝛼1 , 𝛼2 , ⋯ ⋯ , 𝛼𝑛 𝛼𝑖 𝜖 𝐾 .This is the space of row vectors. Addition and Scalar

multiplication are defined by the obvious rules:

𝛼1 , 𝛼2 , ⋯ ⋯ , 𝛼𝑛 + 𝛽1 , 𝛽2 , ⋯ ⋯ , 𝛽𝑛 = 𝛼1 + 𝛽1 , 𝛼2 + 𝛽2 , ⋯ ⋯ , 𝛼𝑛 + 𝛽𝑛 ;

𝜇 𝛼1 , 𝛼2 , ⋯ ⋯ , 𝛼𝑛 = 𝜇𝛼1 , 𝜇𝛼2 , ⋯ ⋯ , 𝜇𝛼𝑛

The most familiar examples are

ℝ2 = 𝑥, 𝑦 : 𝑥, 𝑦 𝜖ℝ

and

ℝ3 = 𝑥, 𝑦, 𝑧 : 𝑥, 𝑦, 𝑧 𝜖ℝ
Which we can think of geometrically as the points in ordinary 2-and 3-dimensional
space, equipped with a coordinate system.

2: Let 𝐾 𝑥 be the set of polynomials in an indeterminate x with coefficients in the field

K.i.e 𝐾 𝑥 = 𝑎0 + 𝑎1 𝑥 + ⋯ ⋯ + 𝑎𝑛 𝑥 𝑛 : 𝑛 ≥ 0, 𝑎𝑖 𝜖 𝐾

Then 𝐾 𝑥 is a vector space over 𝐾.

3: Let 𝐾 𝑥 ≤ 𝑛 be the set of polynomials over K of degree at most n, for some 𝑛 ≥ 0

Then 𝐾 𝑥 ≤ 𝑛 is also a vector space over K; in fact it is a subspace of 𝐾 𝑥 .

Note that the polynomials of degree exactly n do not form a vector space.

4: Let K=ℝ and let V be the set of n-times differentiable functions 𝑓: ℝ → ℝ which are
solutions of the differential equation
𝑑𝑛 𝑓 𝑑𝑛 −1 𝑓 𝑑𝑓
𝜇0 + 𝜇1 + ⋯ ⋯ + 𝜇 𝑛−1 + 𝜇𝑛 𝑓 = 0
𝑑𝑥 𝑛 𝑑𝑥 𝑛 −1 𝑑𝑥

For fixed 𝜇0 , 𝜇1 , ⋯ ⋯ , 𝜇𝑛 𝜖 ℝ .Then 𝑉 is a vector space over ℝ, for if 𝑓(𝑥) and 𝑔(𝑥) are
both solutions of this equation ,then so are 𝑓(𝑥) + 𝑔(𝑥) and 𝛼𝑓 𝑥 𝑓𝑜𝑟 𝑎𝑙𝑙 𝛼 𝜖 ℝ.

5: There are many such examples that are important in Analysis. For example, the set
𝑐 𝑘 0,1 , ℝ , consisting of all functions 𝑓: 0,1 → ℝ such that the kth derivative 𝑓 𝑘
exists and is continuous, is a vector space over ℝ with the usual pointwise definitions of
addition and scalar multiplication of functions.

We shall be assuming the following additional simple properties of vectors and

scalars from now on. They can all be deduced from the axioms.

(i) 𝛼0 = 𝑜 𝑓𝑜𝑟 𝑎𝑙𝑙 𝛼 𝜖 𝐾

(ii) 0v=0 for all v 𝜖 V
(iii) − 𝛼𝑣 = −𝛼 𝑣 = 𝛼 −𝑣 𝑓𝑜𝑟 𝑎𝑙𝑙 𝛼 𝜖 𝐾 𝑎𝑛𝑑 v 𝜖 V
(iv) if 𝛼𝑣 = 0 𝑡𝑕𝑒𝑛 𝛼 = 0 𝑜𝑟 𝑣 = 0
VECTOR SUBSPACE SPANNED BY A GIVEN SYSTEM OF VECTOR: Let 𝑎, 𝑏, 𝑐 be any three
vectors of 𝑉3 . The set 𝑆 of all vectors of the form 𝑥𝑎, 𝑦𝑏, 𝑧𝑐, where 𝑥, 𝑦, 𝑧 are any scalars,
is a subspace of 𝑉3 . If 𝑥1 𝑎 + 𝑦1 𝑏 + 𝑧1 𝑐 be any two members of 𝑆, then 𝑥1 𝑎 + 𝑦1 𝑏 +
𝑧1 𝑐 + 𝑥2 𝑎 + 𝑦2 𝑏 + 𝑧2 𝑐 = 𝑥1 + 𝑥2 𝑎 + 𝑦1 + 𝑦2 𝑏 + 𝑧1 + 𝑧2 𝑐, which is also a
member of 𝑠. Also if 𝑘 be any scaler, then 𝑘 𝑥1 𝑎 + 𝑦1 𝑏 + 𝑧1 𝑐 = 𝑘𝑥1 𝑎 + 𝑘𝑦1 𝑏 +
𝑘𝑧1 𝑐, which is also a member of 𝑆. Also if 𝑘 be any scalar, then 𝑘 𝑥1 𝑎 + 𝑦1 𝑏 + 𝑧1 𝑐 =
𝑘𝑥1 𝑎 + 𝑘𝑦1 𝑏 + 𝑘𝑧1 𝑐, which is again a member of 𝑆. Thus 𝑆 is a vector subspace
and we say that 𝑆 is spanned by the vectors 𝑎, 𝑏 and 𝑐. More generally;

If 𝑎1 , 𝑎2 , … , 𝑎𝑟 be a set of 𝑟 fixed vectors of 𝑉𝑛 , then the set 𝑆 of all 𝑛- vectors of the form
𝑝1 𝑎1 + 𝑝2 𝑎2 + ⋯ + 𝑝𝑟 𝑎𝑟 where 𝑝1 , 𝑝2 , … , 𝑝𝑟 are any scalers is a vector subspace of 𝑉𝑛 .

This vector space is said to be spanned by the vectors 𝑎1 , 𝑎2 , … , 𝑎𝑟 . Thus a vector space
which arises a set of all linear combinations of any given set of vectors is said to be
spanned by the given set of vectors.

VECTOR TRIPLE PRODUCT: The vector product of two vectors one of which is itself the
vector product of two vectors is a vector quantity called a “Vector Triple product”. Thus
if 𝒂, 𝒃, and 𝒄 be three vectors, the product of the form 𝒂 × 𝒃 × 𝒄 and 𝒂 × 𝒃 × 𝒄 etc.
are called “Vector Triple Products”.

VELOCITY POTENTIAL: If 𝑞 is the fluid velocity at any instant 𝑡 then the equation of the
streamline, at that instant are defined as

𝑑𝑥 𝑑𝑦 𝑑𝑧
= =
𝑢 𝑣 𝑤

These curves cut the surfaces

𝑢 𝑑𝑥 + 𝑣 𝑑𝑦 + 𝑤 𝑑𝑧 = 0

orthogonally. Consider a scalar function ∅(𝑥, 𝑦, 𝑧, 𝑡) at the instant, uniform throughout

the entire flow field, such that

−𝑑∅ = 𝑢 𝑑𝑥 + 𝑣 𝑑𝑦 + 𝑤 𝑑𝑧.

∂∅ ∂∅ ∂∅
Or − ∂x dx − ∂y dy − ∂z dz = 𝑢 𝑑𝑥 + 𝑣 𝑑𝑦 + 𝑤 𝑑𝑧.

The expression on the R.H.S. of (2) is exact differential so

∂∅ ∂∅ ∂∅
𝑢= , v = − ∂y , w = − ∂z
∂x

Or q = −∆∅ ,
where ∅ is termed the velocity potential for the filed 𝑞, the negative sign is taken as a
matter of convention.

VENN DIAGRAM:A Venn diagram is a picture that illustrates the relationships between
sets. The universal set you are considering is represented by a rectangle, and sets are
represented by circles or ellipses. The possible relationships between two sets A and B
are as follows:
 Set B is a subset of set A.
 Set A is a subset of set B.
 Set A and set B are disjoint (they have no elements in common).
 Set A and set B have some elements in common.
Venn diagram for the universal set of complex numbers is

VERTEX METHOD: For a linear programming problem where the decision variables are
not required to take integer values, the optimal solutions will occur at one or more of
the extreme points, i.e. vertices of the feasible region. The method of solution is to
calculate the value of the objective function at each vertex, and the maximum or
minimum (as required) of these values identifies any optimal solutions. Where more
than one such vertex is found, all points on the boundary of the feasible region between
them will also be an optimal solution.
VERTICAL LINE: The vertical line, also called the vertical slash or upright slash (|), is
used in mathematical notation in place of the expression "such that" or "it is true that."
This symbol is commonly encountered in statements involving logic and sets.
VERTICAL LINE TEST: The vertical line test can be used to determine if a relation is a
function. If a vertical line can be drawn that crosses two points on the graph of the
relation, then the relation is not a function.
VISCOSITY: Each element of fluid experiences stress exerted on it by other elements of
the fluid which surround it. The stress at each part of the surface of element is resolved
into two components: normal and tangential to the surface, which are called pressure
and shear stress respectively. Pressure is exerted whether the fluid is moving or at rest,
but shear stress occurs only in moving fluids. The feature is the one which enables fluids
to be distinguished from solids. Matter in solid form can exert shear stresses when at
rest, liquids and gases can not. The property which gives rise to shear stresses is called
viscosity. It is a measure of its resistance to flow. It is possessed by all real fluids and its
magnitude is expressed by a coefficient which relates the size of the shear stress at a
point in a fluid to the rate of shear strain which causes it.

The ratio of coefficient of viscosity μ to the mass density ρ of the fluid is known as the
coefficient of kinematic viscosity which is represented as 𝑣 = μ/ρ. The viscosity of a
fluid is practically independent of pressure and depends upon temperature only. The
kinematic viscosity of liquid at a given pressure is a function of temperature. The
dimensions of viscosity can be determined as

τ Shearingstress
μ= =
du/dy Velocityqradient

force/unit area (ML/T 2 ) × (1/L2 )

= = = ML−1 T −1
rate of shear (L/1) × (1/L)
.

The kinematic viscosity, which is the ratio of absolute viscosity to the density, has
dimensions of length and time.

viscosity μ M/LT 2 −1
v= = = LT
density ρ M/L3

Viscosity of the fluid is practically independent of pressure and the depends upon the
temperature only.
VOGEL’S APPROXIMATION MEHTOD OR VAM METHOD OR PENALTY METHOD OR LEAST
ERROR METHOD: For each row and column remaining under consideration, calculate
its difference, the arithmetic difference between the smallest and next-to-the-smallest
unit cost 𝑐𝑖𝑗 still remaining in that row or column. In that row or column having the
largest difference, select the variable having the smallest remaining unit cost. VAM is the
most appropriate and lest error method, widely used for the solution of transportation
problems. The difference steps used for this method can be made well clear by the
following example:

Let us have the following problem to be solved by VAM method.

2 7 4
2 3 3 1
5 4 7
1 6 2

We can solve the above transportation problem by applying VAM

5 2 7 4

3 3 8 1

5 7 4 8 7

2 1 2 6 10 2

7 9 18

2 7 4

3 3 8 1

5 4 7

1 6 2

Penalties 1 1 1
 5 2 7 4

5 4 7

1 6 2

Penalties 1 2 2

5 4 7
2
1 6 10 2

Penalties 4 2 5

5 7 4

2 1 2 6
So, solution set is

5 0 0
0 0 8
0 7 0
2 2 10

Min cost 5 × 2 + 2 × 1 + 7 × 4 + 6 × 2 + 8 × 1 + 10 × 2

= 10 + 2 + 28 + 12 + 8 + 20

= 80

VOLUME: The volume of a solid is a measure of how much space it occupies. The volume
of a cube with edge a units long is 𝑎3 . Volumes of other solids are measured in cubic
units. The volume of a prism or cylinder is (base area) (altitude), and the volume of a
pyramid or cone is (1/3) (base area) (altitude).
VOLUME OF TETRAHEDRON: The volume of a tetrahedron, whose three coterminous
𝟏
edges in the right-handed orientation are 𝒂, 𝒃, 𝒄 is 𝟔 𝒂, 𝒃, 𝒄 .

The four points 𝒂, 𝒃, 𝒄, 𝒅 will be coplanar if the volume of the tetrahedron formed by
them is zero,

i.e. 𝒂 − 𝒅, 𝒃 − 𝒅, 𝒄 − 𝒅 = 𝟎 i.e. 𝒂𝒃𝒄 − 𝒂𝒃𝒅 + 𝒂𝒄𝒅 − 𝒃𝒄𝒅 = 𝟎

Also, the colume of the tetrahedrom is

𝑎1 𝑎2 𝑎3 1
𝑏1 𝑏2 𝑏3 1
1 𝑐1 𝑐2 𝑐3 1
6 𝑑1 𝑑2 𝑑3 1

where 𝑎1 , 𝑎2 , 𝑎3 , 𝑏1 , 𝑏2 , 𝑏3 , 𝑐1 , 𝑐2 , 𝑐3 , 𝑑1 , 𝑑2 , 𝑑3 are the co-ordinates of the vertices.

The volume of a tetrahedron bounded by the four planes 𝒓 ∙ 𝒎𝒋 + 𝒏𝒌 = 𝟎, 𝒓 ∙
𝒏𝒌 + 𝒍𝒊 = 𝟎, 𝒓 ∙ 𝒍𝒊 + 𝒎𝒋 = 𝟎,

2𝑝 3
and 𝒓 ∙ 𝒍𝒊 + 𝒎𝒋 + 𝒎𝒌 = 𝒑 𝑖𝑠 .
3𝑙𝑚𝑛

VON MANGOLDT FUNCTION: We define the Von Mangoldt function 𝛬 ∶ 𝑁 → 𝐶, defined

for every 𝑛 ∇ 𝑁 by writing 𝛬(𝑛) = 𝑙𝑜𝑔 𝑝, if 𝑛 = 𝑝𝑢, with 𝑝 prime and 𝑢 ∇ 𝑁, 0,
otherwise. For every 𝑛 ∇ 𝑁, we have 𝑚 |𝑛 𝛬(𝑚) = 𝑙𝑜𝑔 𝑛.

VON NEUMANN'S TRACE INEQUALITY: It states that for any 𝑛 × 𝑛 complex

matrices 𝐴, 𝐵 with singular values

and respectively, we have

The equality is achieved when and are simultaneously unitarily diagonalizable.

VORTEX LINE: A vortex lien is a line whose direction coincides with the direction of the
instantaneous axis of molecular rotation. In other words, a line to which vorticity vector
are tangent at all its points is called a vortex line. The differential equation to the vortex
line is

dx dy dz
ω × dr = 0 ⇒ = =
ξ η ζ

VORTEX TUBE: A vortex tube is obtained if through every point of a small closed curve
(tube), the corresponding vortex cannot originate or terminate at internal points in
fluid. They can only form closed or terminate on boundaries. In the case of smoke rings
the vortex lines form closed curves where as in a whirlpool it terminate on the
boundary of the fluid.

 The fluid contained within the vortex tube constitutes the vortex filaments or
simply vortices. The boundary of a vortex filament is called a vortex tube.
 Every vortex is always composed of the same element of fluid.
 The product of the angular velocity of any vortex into its cross-section is
constant with respect to the time, and is the same throughout its length.
 Every vortex must either form a closed curve or have its extrenities in the
boundaries of fluid.

VORTICITY VECTOR: Let 𝑞 represents the velocity of a fluid motion, then the vorticity of
an element of fluid is defined as the curl of its velocity vector. That is, vorticity ω is
defined by ω = curl q.

Physically, vorticity may be generated in an inviscid fluid by the rotational body forces.

WAITING TIME: It is the time up to which a unit has to wait before it is taken into
service after arriving at the servicing station. This a studied with the help of waiting
time distribution.

The waiting time depends on

(a) The number of units already there in the system,

(b) The number of servicing stations in the system,
(c) The schedule in which units are selected for service.
WALK (GRAPH THEORY): The general name given to a sequence of edges and vertices
in a graph. Walks with certain properties are of particular interest and are given specific
names. So a walk in which no edge is repeated is a trail, and if no vertex is revisited in
the course of a trail it is a path. The special case where the path finishes at the vertex
where it started is an exception to this, and is termed a cycle or loop.

WALSH-HADAMARD CODE: The binary Walsh-Hadamard code of dimension 𝑘 is defined

by the following function 𝐶. For every 𝑥 ∇ 𝐹2𝑘 , 𝐶 𝑥 = (⌌𝑥, 𝑧1 ⌍, . . . , ⌌𝑥, 𝑧2𝑘−1 ⌍), where 𝑧1 ,
𝑘
. . . , 𝑧2𝑘−1 are all the nonzero words in 𝐹2𝑘 , and ⌌𝑥, 𝑧⌍ ≡ 𝑖=1 𝑥𝑖 · 𝑧𝑖 (𝑚𝑜𝑑 2).
WARING’S PROBLEM: The first formulation of Waring’s problem is found in E. Waring,
“Meditutiones algebraicae” (1770) in which he discusses the problem of expressing an
arbitrary positive integer as the sum of at most nine cubes or as the sum of at most
nineteen biquadratics.
WAVE MOTION: Wave motion is due to the action of gravity which acts in the direction
of restoring the undisturbed state of rest. Wave motion of a liquid acted upon by gravity
and having a free surface is a motion in which the elevation of the free surface above a
fixed horizontal plane varies.

A simple harmonic progressive wave represented by a sine curve moving with definite
velocity in the direction of its length is of the form

y = a sin mx − nt = a sin m x − (n/m)t

which shows that the profile y = a sin mx at t = 0 move, with velocity c(= n/m) along
the positive direction of the X- axis, c is called the velocity of propagation of the wave.

The maximum value of the disturbance y, viz. , a is called the amplitude of the wave. The
points P and P′ of maximum elevation are called the troughs of the wave. The distance
between successive crests is called the wave length and is denoted by λ i. e. , λ = 2π/m.

The aspect of the free surface is same at time t and t + 2π/n. Thus the period, T of a
wave is 2π/n. The reciprocal of a period is called the frequency. It denotes the number
of oscillations per second.

WEAK ARROW- HURWICZ- UZAWA CONSTRAINT: Let X 0 be an open set in Rn and let 𝑔
be an 𝑚- dimensional vector function defined on X 0 , and let

𝑋 = x: x ∇ X 0 , g(x) ≤ 0 .
𝑔 is said to satisfy the Weak Arrow- Hurwicz- Uzawa constraint qualification at x ∇ X if
𝑔 is differentiable at x, and ∆g Q x z > 0, ∆g Q x z ≥ 0 has a solution z ∇ Rn , where
𝑃 = i: g i x = 0, and g i is pseudoconcave at x and
𝑄 = i ∶ g i x = 0, and g i is pseudoconcave at x

WEAK CONVERGENCE: Let 𝑁 be a normed linear space. The sequence 𝑓𝑛 of 𝑁 is said to

𝜔
converge weakly to 𝑓 ∇ 𝑁, written as 𝑓𝑛 → 𝑓, if for every Φ 𝜖 𝑁 ∗ (the dual space of 𝑁),

Φ 𝑓𝑛 ⟶ Φ 𝑓 , 𝑎𝑠 𝑛 ⟶ ∞

Here 𝑓 is said to be weak limit of the sequence 𝑓𝑛

WEAK DUALITY THEOREM FOR LP: If 𝑥 is feasible for 𝑃 (so 𝐴𝑥 ≤ 𝑏, 𝑥 ≥ 0) and 𝜆 is

feasible for 𝐷 (so 𝜆 ≥ 0, 𝐴𝑇 𝜆 ≥ 𝑐) then 𝑐 𝑇 𝑥 ≤ 𝜆𝑇 𝑏.

WEAK DUALITY THEOREM: Let X 0 be open, and let θ and g be differentiable on X 0 . Then

x1 ∇ X, x 2 , u2 ∇ Y ⇒ θ x1 ≥ Ψ x 2 , u2

where θ and g are convex at x 2 and X and Y are defined in the (primal) minimization
problem and the dual (maximization) problem.

WEAK REVERSE CONVEX CONSTRAINT QUALIFICATION: Let X 0 be an open set in Rn and

let 𝑔 be an 𝑚- dimensional vector function defined on X 0 , and let

𝑋 = x: x ∇ X 0 , g(x) ≤ 0 .

𝑔 is said to satisfy the Weak reverse convex constraint qualification at x ∇ X if 𝑔 is

differentiable at x, and for each i ∇ I, either g1 pseudoconave at x or g i is linear on Rn ,
where 𝐼 = i ∶ g i x = 0

WEAK TOPOLOGY AND STRONG TOPOLOGY: Let 𝑋 be a normed linear space and 𝑋’ its
dual space. Take a finite number of elements 𝑥1′ , 𝑥2′ , − − −, 𝑥𝑛′ from 𝑋’, and consider the
subset of X: {𝑥 ∇ 𝑋; 𝑠𝑢𝑝 ⌌𝒙, 𝒙′ ⌍ ≤∇}, ∇> 𝑂. If we take the totality of such subsets of 𝑋 as a
fundamental system of neighborhoods of 0 of 𝑋, then 𝑋 is a locally convex topological
linear space, denoted sometimes by 𝑋𝑤 . This topology is called the weak topology of 𝑋. If
a sequence {𝑥𝑛 } ⊆ 𝑋 converges to 𝑥 ∇ 𝑋 with respect to the weak topology of 𝑋, then
it is said to converge weakly. This is equivalent to the convergence (𝑥𝑛 , 𝒇 ) → (𝒙, 𝒇) for
any 𝒇 ∇ 𝑋’. The original topology of 𝑋 determined by the norm is then called the strong
topology of 𝑋, and to stress the strong topology we sometimes write 𝑋𝑠 in place of the
original 𝑋.
WEIERSTRASS APPROXIMATION THEOREM: For any function 𝑓 ∇ 𝐶[𝑎, 𝑏] and any є > 0
there exists a polynomial 𝑝 such that ||𝑓 − 𝑝||∞ < є.

WEIERSTRASS M-TEST: If 𝑢𝑛 (𝑧) ≤ 𝑀𝑛 where 𝑀𝑛 is independent of 𝑧 in domain 𝑅

and ΣM𝑛 , the series of positive constant is convergent, then the series Σ𝑢𝑛 (𝑧) is
uniformly convergent.

WEIERSTRASS, KARL (1815–97): German mathematician who was a leading figure in

the field of mathematical analysis. His work was concerned with providing the subject
with the necessary rigour as it developed out of eighteenth-century calculus. One
particular area in which he made significant contributions was in the expansion of
functions in power series. To show that intuition is not always reliable, he gave an
example of a function that is continuous at every point but not differentiable at any
point. Some of his work was done while he was a provincial school teacher, having little
contact with the world of professional mathematicians. He was promoted directly to
professor of mathematics in Berlin at the age of 40.
WEIERSTRASS’S ELLIPTIC FUNCTION: Weierstrass function p(z) is defined by the
equation
d
p z =− ξ z .
dz
∞
1 1 1
⇒p z = 2+ 2
,
z −Ωm, n (Ω2 m, n)
m,n − −∞

where summation is extended over all positive and negative integral and zero values of
𝑚 and 𝑛, save 𝑚 = 𝑛 = 0.

p(z) is an analytic function whose only singularities are double poles of residue zero at
each of the point Ω𝑚 ,𝑛 .

WEIERSTRESS’S SIGMA FUNCTIONS: Weierstrass’s Sigma function is given by

∞
𝑧 𝑧 𝑧
ς 2/ω1 , ω2 = z 1− 𝑒𝑥𝑝. +
Ωm,n Ωm,n 2Ωm,n
𝑚 ,𝑛=∞
where multiplication is extended over all positive and negative integral and zero values
of 𝑚 = 𝑛 = 0.

ς 2/ω1 , ω2 is an integral function of order 2, with simple zeroes at the points Ωm,n .
This function is also denoted by ς(z).

WEIERSTRASS’S THEOREM COMPLEX ANALYSIS): If 𝑧1 , 𝑧2, … . , 𝑧𝑛 be any sequence of

numbers whose only limiting point is the point at infinity, then it is possible to construct
an integral function which vanishes at each of these points 𝑧𝑛 .

WEIERSTRASS’S THEOREM (ESSENTIAL SINGULARITY): If 𝑧 = 𝛼 is an essential

singularity of 𝑓(𝑧), given any positive numbers, 𝑟, 𝜀 and any number 𝑐, there is a point in
the circle 𝑧 − 𝛼 < 𝑟 at which 𝑓 𝑧 − 𝑐 < 𝜀.

Or, In other words, in every arbitrary neighborhood of an essential singularity, there

exists a point (and therefore an infinite number of points) at which the function differs
as little as we place from any pre-assigned number.

WEIERSTRASS THEOREM: For any finite 𝐼 = [𝑎, 𝑏], 𝑃 is dense in 𝐶(𝐼), i.e., for each
𝑓 ∇ 𝐶(𝐼) and for each 𝜀 > 0 there exists some 𝑝 ∇ 𝑃 such that |𝑓(𝑥) − 𝑝(𝑥)| <
𝜀, 𝑎 ≤ 𝑥 ≤ 𝑏 .
WELL-FORMED FORMULA: A well-formed formula (or wff ) is a sequence of symbols
that is an acceptable formula in logic. For example, the sequence p AND q is a wff, but
the sequence AND pq is not a wff. Certain rules govern the formation of wff ’s in a
particular type of logic. Here is an example of such a rule: If p and q are wff’s, then (p
AND q) is also a wff.
WELL-ORDERING THEOREM: Every nonempty set 𝑆 of nonnegative integers contains a
least element; that is, there is some integer 𝑎 in 𝑆 such that 𝑎 ∇ 𝑏 for all 𝑏 belonging to
𝑆. In 1904, E. Zermelo first stated the axiom of choice and used it for his proof of the well
ordering theorem, which says that every set can be well-ordered by an appropriate
ordering. Conversely, the well-ordering theorem implies the axiom of choice. Many
important results in set theory can be obtained by using the axiom of choice, for
example, that Cardinal numbers are comparable, or that various definitions of the
finiteness or infiniteness of sets are equivalent. Various important theorems outside of
set theory, e.g., the existence of bases in a linear space, compactness of the direct
product of compact topological spaces Tikonov’s theorem , the existence of a subset
which is not Lebesgue measurable in Euclidean space, etc., are proved using the axiom
of choice.

WEYL'S INEQUALITY: If 𝑀, 𝑁, 𝑎 and 𝑞 are integers, with 𝑎 and 𝑞 coprime, 𝑞 > 0,

and 𝑓 is a real polynomial of degree 𝑘 whose leading coefficient 𝑐 satisfies

for some 𝑡 greater than or equal to 1, then for any positive real number ∇ one has

This inequality will only be useful when 𝑞 < 𝑁 𝑘 .

WHITNEY SUM: A Whitney sum is an analog of the direct product for vector bundles.
Given two vector bundles 𝛼 and 𝛽 over the same base 𝐵 their Cartesian product is a
vector bundle over 𝐵 × 𝐵. The diagonal map induces a vector bundle
over 𝐵 called the Whitney sum of these vector bundles and denoted by 𝛼 ⊕ 𝛽.

WHITNEY THEOREM: Let 𝑀 be an abstract smooth manifold of dimension 𝑚, and

assume there exists a countable atlas for 𝑀. Then there exists a diffeomorphism of 𝑀
onto a manifold in 𝑅 2𝑚 .
WILES, ANDREW JOHN FRS (1953–): British mathematician famous for the proof of one
of mathematics most famous problems—Fermat’s Last Theorem. He announced he had
the proof in 1993 and was widely tipped to receive a Fields Medal as a result. However,
that work was found to contain an error and it was not until 1995 that a complete proof
was published, by which time Wiles was too old to be considered for a Fields Medal.
However, the Fields Institute honoured his extraordinary achievement by a special
award in 1998 of a silver plate. He also won the Wolf Prize in 1996.
WINDING NUMBER: A cycle 𝛤 consists of finitely many closed piecewise smooth
contours 𝛾1 , . . . . . , 𝛾𝑛 , and the integral over 𝛤 is the sum of the integrals over the 𝛾𝑘 . If a is
a point not lying on 𝛤 (i.e. not on any of the 𝛾𝑘 ) we define
1 𝑛 . 1
𝑛(𝛤, 𝑎) = 𝑘=1 ∫𝛾 𝑘 𝑧 − 𝑎 𝑑𝑧, which will be called the winding number of 𝛤 about 𝑎.
2𝜋𝑖

We have:
(i) 𝑛(𝛤, 𝑎) is an integer.
 (ii) If 𝜍 ∶ [𝑐, 𝑑] → 𝐶 is a path not intersecting 𝛤, then 𝑛(𝛤, 𝜍(𝑡)) is
constant, and hence if 𝐷 is a domain in 𝐶 not intersecting 𝛤, then
𝑛(𝛤, 𝑎) is constant on 𝐷.
(iii) If |𝑤| is large enough, 𝑛(𝛤, 𝑤) = 0.
(iv) If 𝛤 is the circle |𝑧 − 𝑎| = 𝑠 > 0, described once counter-
clockwise, then 𝑛(𝛤, 𝑎) = 1.
WLOG: “WLOG” or “WOLOG” is an acronym which stands for “without loss of
generality.” WLOG is invoked in situations where some property of a model or system is
invariant under the particular choice of instance attributes, but for the sake of
demonstration, these attributes must be fixed.
WOLFE’S DUALITY THEOREM: Let X 0 be an open set in Rn and let 𝜃 and 𝑔 be
differentiable and convex on X 0 , let x solves maximization problem. Let g satisfy any one
of the six constraint qualifications. Then there exists u ∇ Rm such that (x, u) solves the
dual problem (maximization) DP and θ x = Ψ x, u .

WORK: Force acting on a particle does work when the particle is displaced in a direction
which is not perpendicular to the force. The work done is a scalar quantity and its
measure is equal to the product of the magnitude of the force and the resolved part of
the displacement in the direction of the force. Thus if F, d are vectors representing the
force and displacement respectively inclined at an angle 𝜃, the measure of the work
done is = F. d i. e. w = Fd cos θ = F d cosθ, where F = F and d = d .

From the above, it is clear that the work done is zero only where 𝜃 = 𝜋/2 i.e. when d is
perpendicular to F.

WORK–ENERGY PRINCIPLE: The principle that the change in kinetic energy of a particle
during some time interval is equal to the work done by the total force acting on the
particle during the time interval.
x, y COORDINATES: 𝑥, 𝑦 coordinates are respectively the horizontal and vertical
addresses of any pixel or addressable point on a computer display screen.
The x coordinate is a given number of pixels along the horizontal axis of a display
starting from the pixel (pixel 0) on the extreme left of the screen. The y coordinate is a
given number of pixels along the vertical axis of a display starting from the pixel (pixel
0) at the top of the screen. Together, the x and y coordinates locate any specific pixel
location on the screen. x and y coordinates can also be specified as values relative to any
starting point on the screen or any subset of the screen such as an image. On the Web,
each clickable area of an image map is specified as a pair of x and y coordinates relative
to the upper left-hand corner of the image.
X-INTERCEPT: The 𝑥-intercept of a curve is the value of 𝑥 at the point where the curve
crosses the 𝑥 -axis.
Y-INTERCEPT: The 𝑦-intercept of a curve is the value of 𝑦 at the point where the curve
crosses the 𝑦-axis.
YOUNG’S INEQUALITY: Let two real numbers 1 < 𝑝, 𝑞 < ∞ are related through
𝑎 𝑝 𝑏 𝑞
1/𝑝 + 1/𝑞 = 1 then 𝑎𝑏 ≤ + for any complex 𝑎 and 𝑏.
𝑝 𝑞

YOUNG’S MODULUS OF ELASTICITY: The modulus of elasticity of elastic strings or

springs made of the same material, but of different cross-sectional area will vary and
𝜆
Young’s modulus of elasticity states the relationship as 𝐸 = 𝐴 giving Hooke’s law
𝐸𝐴
𝑇= 𝑥 as where 𝐸 = Young’s modulus of elasticity, 𝑥 = modulus of elasticity, 𝐴 =
𝐿

cross-sectional area, 𝑙 = natural length and 𝑥 = extension from natural length.

ZARISKI TOPOLOGY: The Zariski topology on an algebraic set 𝑉 in 𝐴𝑛 is the topology
whose open sets are of the form 𝑉 \𝑉 (𝐼) for some ideal 𝐼 of 𝐾[𝑋1 , 𝑋2 , . . . , 𝑋𝑛 ].
ZEEMAN, SIR ERIK CHRISTOPHER FRS (1925- ): British mathematician whose work in
topology and catastrophe theory has found wide applicability in areas as diverse as
physics, the social sciences and economics.
ZENO’S PARADOX: Zeno’s paradox claims that an object can never travel a distance 𝑑
because it first must pass through the point 𝑑/2; before that it must pass the point 𝑑/4;
before that it must pass the point 𝑑/8; and so on. Since there are an infinite number of
points, Zeno’s paradox claims that it would take an infinite amount of time. Since in
reality objects can move from one point to another, Zeno’s paradox is based on a
misunderstanding of continuous space. Alternatively, space might not be continuous on
extremely small scales, in which case an object does jump from one location to an
adjacent location without passing through any intermediate locations. In any case, the
laws of quantum mechanics make it impossible to measure the exact location of
something with perfect accuracy.
ZERMELO, ERNST (FRIEDRICH FERDINAND) (1871-1953): German mathematician
considered to be the founder of axiomatic set theory. In 1908, he formulated a set of
axioms for set theory, which attempted to overcome problems such as that posed by
Russell’s paradox. These, with modifications, have been the foundation on which much
subsequent work in the subject has
been built.
ZERMELO-FRAENKEL SET THEORY: ZF is a formal system expressed in ‘first-order
predicate logic with the predicate symbol= (equality) and based on the some axioms.
The axioms do not contain any predicate symbol other than 𝜖, where 𝑥𝜖𝑦 is read “𝑥 is an
element of 𝑦.” Any formula containing only 𝜖 as a predicate symbols is called a set-
theoretic formula. Following the usual convention, we omit the outermost universal
quantifier, and use restricted quantifiers such as ∃𝑥 𝜖 𝑎, 𝑉𝑥 𝜖 𝑎.

1. AXIOM OF EXTENSIONALITY:

∀𝑥 𝑥 𝜖 𝑎 ≡ 𝑥 𝜖 𝑏 → 𝑎 = 𝑏.

This asserts that sets formed by the same elements are equal. The formula 𝑥 𝜖 𝑎 (𝑥 𝜖 𝑏)
is denoted by 𝑎 ⊂ 𝑏. This means “𝑎 is a subset of 𝑏.” Then Axiom 1 can be expressed by

𝑎 ⊂ 𝑏 ∧ 𝑏 ⊂ 𝑎 → 𝑎 = 𝑏.

2. AXIOM OF THE UNORDERED PAIR:

∃𝑥∀𝑦 𝑦 𝜖 𝑥 ≡ 𝑦 𝜖 𝑎 ∨ 𝑦 = 𝑏 .

This asserts the existence, for any sets a and b, of a set 𝑥 having a and b as its only
elements. This 𝑥 is called the unordered pair of a and b and is denoted by 𝑎, 𝑏 . The
notion of ordered pair is characterized by

𝑎, 𝑏 = 𝑐. 𝑑 ≡ 𝑎 = 𝑐 ∨ 𝑏 = 𝑑.

An example of such is 𝑎, 𝑏 = 𝑎, 𝑎 , 𝑎, 𝑏 .

3. AXIOM OF THE SUM SET:

∃𝑥∀𝑦 𝑦 𝜖 𝑥 ≡ ∃𝑧 𝜖 𝑎(𝑦𝜖𝑧) .

This asserts the existence, for any sets a of the sum (or union) 𝑥 of all the sets that are
elements of a. This 𝑥 is denoted by the 𝑈𝑎 or 𝑆(𝑎). We write a 𝑈𝑏 for 𝑈 𝑎, 𝑏 and
𝑎′ for 𝑎𝑈{𝑎. 𝑎}.
 4. AXIOM OF THE POWER SET:

∀𝑥∃𝑦 𝑦 𝜖 𝑥 ≡ ∀𝑧 𝜖 𝑦(𝑧𝜖𝑎) .

This asserts the existence, for any sets a of the power set 𝑥 consisting of all the subsets
of a. This 𝑥 is denoted by 𝑃𝑎. We have 𝑆(𝑃 𝑎 = 𝑎, 𝑠𝑜 𝑆 is a left inverse of 𝑃 and the
products 𝑆𝑃 𝑎𝑛𝑑 𝑃𝑆 are idempotent.

5. AXIOM OF THE EMPTY SET:

∃𝑥∀𝑦 ℸ𝑦 𝜖 𝑥 .

This asserts the existence of the empty set. The empty set is denoted by ⊘ or 0.

6. AXIOM OF INFINITY:

∃𝑥 0 𝜖 𝑥 ∨ ∀ 𝑦 𝜖 𝑥(𝑦 ′ 𝜖𝑥) .

This asserts the existence of the set consisting all the natural numbers, whereas
0,1 = 0′ = 0 , 2 = 1′ = 0,1 , 3 = 2′ = 0,1,2 . This definition of natural number is due
to von Neumann.

7. AXIOM OF SEPARATION:

∃𝑥∀𝑦 𝑦 𝜖 𝑥 ≡ 𝑦 𝜖 𝑎 ⋀ 𝐴 (𝑦) .

This asserts that the existence for any set a and a formula 𝐴(𝑦) of a set 𝑥 consisting of all
element of a satisfying 𝐴(𝑦). This 𝑥 is denoted by {𝑦𝜖𝑎 𝐴 𝑦 }. This is rather a schema for
an infinite number of axioms, for there are an infinite number of 𝐴(𝑦). This axiom, also
called the axiom of comprehension or axiom of subsets, was introduced by Zermelo.

For example, the set of all natural numbers is introduced by

{𝑥𝜖𝑎}∀𝑦(0𝜖𝑦 ⋀ 𝐴 𝑧𝜖𝑦 (𝑧 ′ 𝜖𝑦) → 𝑥𝜖𝑦)}.

Where 𝑎 is a satisfying Axiom 6. The set of all natural numbers is denoted by 𝜔 or 𝑁.

8. AXIOM OF REPLACEMENT:

∃𝑥∀𝑦𝜖𝑎 ∃𝑧 𝐴(𝑦, 𝑧) → ∃𝑧 𝜖 𝑥𝐴 (𝑦, 𝑧) .

This asserts the existence for any set a of a set 𝑥 such that for any 𝑦 of 𝑎𝑖 , if there exists a
𝑧 satisfying 𝐴(𝑦, 𝑧) then such 𝑧 exists in 𝑥. If the relation 𝐴(𝑦, 𝑧) determines a function,
then the image of a set by the relation is included in a set, so by Axiom 7, the image is a
set. This axiom was introduced by Frankel.

9. AXIOM OF REGULARITY:

∀𝑥(∀𝑦𝜖𝑥𝐴(𝑦) → ∀(𝑥) → ∀ 𝑥𝐴 (𝑥).

Using this we can show that ℸ𝑥𝜖𝑥,

ℸ(𝑥𝜖𝑦 𝐴 𝑦𝜖 𝑥), etc. If we assume the axiom of choice stated below, then this is equivalent
to the nonexistence of an infinite descending sequence

𝑥𝑛 ∇ ⋯ ∇ 𝑥2 ∇ 𝑥1 .

If a model e1 of a set theory satisfies the axiom of regularity and has an infinite
descending sequence that is not in the model, then such a model is called a nonstandard
model.

10. AXIOM OF CHOICE:

∀𝑥 ∇ 𝑎∃𝑦𝐴(𝑥, 𝑦) → ∃𝑦∀𝑥 ∇ 𝑎𝐴(𝑥, 𝑦 𝑥 ).

This asserts that if for any element 𝑥 of a there is a set 𝑦 such that 𝐴(𝑥, 𝑦), then there is a
choice function 𝑦 for the formula, i.e., 𝐴 𝑥, 𝑦 𝑥 for all 𝑥 in a. Usually a function is
represented by its graph. A set 𝑓 is called a function defined on a if

∀𝑥∀𝑦(𝑥, 𝑦) ∇ 𝑓 → 𝑥 ∇ 𝑎), ∀𝑥 ∇ 𝑎∃𝑦(𝑥, 𝑦) ∇ 𝑓).

∀𝑥 ∇ 𝑎∀𝑦𝐴𝑧 𝑥, 𝑦 ∇ 𝑓⋀ 𝑥, 𝑧 ∇ 𝑓 → 𝑦 = 𝑧).

This formula is denoted by Fnc(𝑓); then the formula 𝐴 𝑥, 𝑓 𝑥 is an abbreviation of

Fun 𝑓 ⋀ 3𝑦(𝑥, 𝑦) ∇ 𝑓 ⋀ 𝐴 (𝑥, 𝑦)).

The axiom of choice is equivalent to many properties, such as the well-ordering

theorem, Zorn’s lemma, and Tikhonov’s theorem on the product of compact space, and
it is used widely and essentially 1 to 9 is called Zermelo-Fraenkel set theory and is
denoted by ZF; the system ZF minus the axiom of replacement is called Zermelo set
theory, denoted by Z; and the system ZF plus the axiom of choice is denoted by ZFC.
The system Z is weaker then ZF, indeed, the existence of the set 𝜔 of all the natural
numbers and of 𝑃 𝑟𝑜 , 𝑃(𝑃 𝜔 ), can be proved in Z, but the existence of the set
𝜔, 𝑝 𝜔 , 𝑃(𝑃 𝜔 , . cannot be proved in Z. However, we can prove its existence in ZF.

The theory ZF minus the axiom of infinity is called general set theory.

ZERO: Intuitively, zero means nothing—for example, the score that each team has at the
beginning of a game is zero. Formally, zero is the identity element for addition, which
means that, if you add zero to any number, the number remains unchanged. In our
number system the symbol “0” also serves as a placeholder in the decimal
representation of a number. Without zero we would have trouble telling the difference
between 1000 and 10.
ZERO DIVISORS: Non-zero elements in a ring whose product is zero. There are no zero
divisors within the real or complex numbers, but there are in other systems. For
1 0 0 0 0 0
example, if 𝑨 = and 𝑩 = then AB=BA= , so 𝑨 and 𝑩 are zero
0 0 0 1 0 0
divisors within the ring of 2 × 2 matrices, defined with the usual matrix addition and
multiplication.
ZERO ELEMENT: An element 𝑧 is a zero element for a binary operation ° on a set 𝑆 if, for
all 𝑎 in 𝑆, 𝑎 ° 𝑧 = 𝑧 ° 𝑎 = 𝑧. Thus the real number 0 is a zero element for multiplication
since, for all 𝑎, 𝑎° = °𝑎 = 0. The term ‘zero element’, also denoted by 0, may be used
for an element such that 𝑎 + 0 = 0 + 𝑎 = 𝑎 for all 𝑎 in 𝑆, when S is a set with a
binary operation + called addition. Strictly speaking, this is a neutral element for the
operation +.
ZERO FUNCTION: In real analysis, the zero function is the real function 𝑓 such that
𝑓(𝑥) = 0 for all 𝑥 in 𝑹.
ZERO MATRIX: The 𝑚 × 𝑛 zero matrix 𝑶 is the 𝑚 × 𝑛 matrix with all its entries zero. A
zero column matrix or row matrix may be denoted by 𝟎.
ZERO OF AN ANALYTIC FUNCTION: A zero of an analytic function 𝑓(𝑧) is a value of 𝑧
such that 𝑓 𝑧 = 0.

An analytic function 𝑓 𝑧 is said to have a zero of order 𝑚 if 𝑓 𝑧 is expressible as

𝑓 𝑧 = (𝑧 − 𝑎)𝑚 ∅ 𝑧 ,
Where ∅(𝑧) is analytic and ∅ 𝛼 ≠ 0. 𝑓(𝑧) is said to have a simple zero at 𝑧 = 𝛼 if 𝑧 = 𝛼
is a zero of order one.

𝑛
If there exists no finite value of 𝑛 𝑠. 𝑡. lim𝑧→𝛼 𝑧 − 𝛼 𝑓 𝑧 = 𝑐 = finite non-zero
constant, then 𝑧 = 𝛼 is called essential singlularity.

ZERO OR NULL VECTOR: The zero or the null vector is a vector whose modulus is zero,
and whose direction is indeterminate. The null vector is represented by the symbol 0. In
the case of the null vector, the initial and terminal points coincide. Thus 𝐴𝐴, 𝑂𝑂 etc. are
null vectors.

ZERO-ONE LAWS: In probability theory there are many theorems claiming that an event
with certain properties has probability 0 or 1. Such theorems are called zero-one laws,
we mention two famous examples, Kolmogorov’z zero-one law and the Hewitt- Savage
zero-one law. Let 𝛼 = 𝛼(𝑋1 , 𝑋2 , … ) be an event concerning a sequence of random
variables 𝑋𝑛 if for every 𝑛, occurrence or nonoccurrence of 𝛼 depends only on
𝑋𝑛 , 𝑋𝑛+1 , … . For example, lim𝑛→∞ 𝑋𝑛 = 0 is a tail event, 𝛼 is called a symmetric event
concerning 𝑋𝑛 if occurrence or nonoccurrence of 𝛼 is invariant under every finite
𝑛
permutation of 𝑋1 , 𝑋2 , ⋯ For example, the event that 𝑘=1 𝑋𝑘 > 0 for infinitely many 𝑛’s
is a symmetric event. Kolmogorov’s zero-one law. Every tail event concerning a
sequence of independent random variables has probability 0 or 1.

ZERO SUM AND NON-ZERO SUM GAME: A game is known as zero sum game in which the
net gain or not profit after the game is zero means nothing comes from outside,
payment is always between the players , the loss of one player is the gain of others , and
a game which is not zero sum game is known as non- zero sum game.

ZORN’S LEMMA: An ordered set 𝑋 is called an inductively ordered set if every totally
ordered subset of 𝑋 has an upper bound. A condition 𝐶 for sets is called a condition of
finite character if a set 𝑋 satisfies C if and only if every finite subset of 𝑋 satisfies 𝐶. A
condition 𝐶 for functions is called a condition of finite character if 𝐶 is a condition of
finite character for the graph of the function. Zorn’s lemma can be stated in any one of
the following ways, which are all equivalent to the axiom of choice. It is often more
convenient to use than the axiom of choice or the well-ordering theorem.
(1) Every inductively ordered set has at least one maximal element.
(2) If every well-ordered subset of an ordered set 𝑀 has an Upper bound, then there is
at least one maximal element in 𝑀.
(3) Every ordered set 𝐴 has a well-ordered subset W such that every Upper bound of 𝑀
belongs to 𝑊 .
(4) For a condition 𝐶 of finite character for sets, every set 𝑋 has a maximal (for the
relation of the inclusion) subset of 𝑋 that satisfies 𝐶.
(5) Let 𝐶 be a condition of finite character for functions from 𝑋 to 𝑌. Then, in the set of
functions that satisfy 𝐶, there is a function whose domain is maximal (for the relation of
the inclusion)