Documenti di Didattica
Documenti di Professioni
Documenti di Cultura
MATEMÁTICAS
EM BREVE RELATO DAS
PRINCIPAIS CONTRIBUIÇÕES
+ m p(1 – p )
1ª EDIÇÃO - 2012
PRÓLOGO
NOSSA GERAÇÃO DEVE SE SENTIR PRIVILEGIADA EM PRESENCIAR O ESTÁGIO ATUAL DO
DESENVOLVIMENTO CIENTÍFICO E TECNOLÓGICO, COM DESTAQUE PARA AS CONQUISTAS
NAS ÁREAS:
2
DE COORDENAÇÂO DE EQUIPES DE PESQUISA, SÂO OS QUE CONSEGUIRAM UMA
MELHOR VISÂO GERAL.
----------------xxx------------------
3
APRESENTAÇÃO..............................................................................6
I- CRONOLOGIA ANTIGA .....................................................................11
II- O SURGIMENTO DA CIÊNCIA GREGA...............................................12
III-CONTRIBUIÇÃO DOS PESQUISADORES
ÁRABES/ISLÂMICOS..........................................................................27
EVOLUÇÃO À PARTIR DO MÉTODO CIENTÍFICO
IV-PRINCÍPIOS DA ÓTICA GEOMÉTRICA, LEI DE HOOKE e
AS LEIS DA MECÂNICA CLÁSSICA.......................................................31
CÁLCULO DIFERENCIAL E INTEGRAL..................................................38
EQUAÇÕES DIFERENCIAIS ORDINÁRIAS.............................................41
EXEMPLOS DE FORÇAS: GRAVIDADE (37), ATRITOS (42), RESISTÊNCIA
DO AR (44), ELÉTRICA e MAGNÉTICA (45)
4
A EQUAÇÃO GERAL, EQUAÇÕES INTEGRAIS.................................120/1
PROBABILIDADE e ESTATÍSTICA………………………….......................…...246
CONTROL THEORY...........................................................................271
PERSPECTIVES……………………………………………………………………………….291
OBSERVAÇÕES…………………………………………………………...................…294
5
APRESENTAÇÃO
- DESDE A FAMOSA FRASE DE SÓCRATES (SEC V AC): “O ARTÍFICE SABE MAIS DO SEU OFÍCIO DO
QUE O FILÓSOFO (DA ÉPOCA) DE SUA FILOSOFIA”, VEM SENDO CONSTATADO PELOS INTELECTUAIS
(E CONFIRMADO PELAS LIMITAÇÕES DE SEUS ENTENDIMENTOS E TEORIAS), E TAMBÉM PELOS
CIENTISTAS, QUE, EM GERAL, A CONTEMPLAÇÃO INTELECTUAL CRITERIOSA DE UMA ATIVIDADE SE
CONSTITUI EM TAREFA MAIS LABORIOSA DO QUE A SUA PRÁTICA.
“TODAS AS COISAS PODEM SER REDUZIDAS À SUA ESSÊNCIA NUMÉRICA (ATRAVÉS DOS NÚMEROS
RACIONAIS, OU SEJA, QUOCIENTES DE NÚMEROS INTEIROS)”;
6
- A MECÂNICA DE NEWTON É CONSIDERADA A 2ª GRANDE SÍNTESE MATEMÁTICA, AO, NO ESPAÇO
DA GEOMETRIA EUCLIDIANA, OBTER LEIS DE MOVIMENTO, QUE DESCREVEM TANTO OS
MOVIMENTOS TERRESTRES, QUANTO OS CELESTES, QUE DESDE ARISTÓTELES ERAM TRATADOS
SEPARADAMENTE. ALÉM DISSO, CONSEGUIU PREVER A EXISTÊNCIA DE PLANETAS.
7
NÃO DEVE SER MEMBRO DA COLEÇÃO”, E LEVANTANDO O PROBLEMA DE QUAL SERIA A REGRA DE
FORMAÇÃO DOS CONJUNTOS......
- ATUALMENTE, TEMOS TEORIAS MATEMÁTICAS PRONTAS PARA SEREM APLICADAS, E QUE TEM
APLICAÇÃO (COMO ALGUMAS DAS TEORIAS MENCIONADAS), E TEORIAS EM ELABORAÇÃO, A
MAIORIA INSPIRADA EM PROBLEMAS DA PRÁTICA, QUE, DEPENDENDO DOS RESULTADOS (A SEREM)
OBTIDOS, PODERÃO VIR A TER APLICAÇÃO, TAMBÉM A EXEMPLO DAS TEORIAS MENCIONADAS.
8
SER “VALIDADOS APENAS PELO QUESTIONAMENTO TEÓRICO” , TENDO PREVALECIDO, O MÉTODO
CIENTÍFICO DE ALHAZEN/ GALILEU.
- EM 1781, EM SUA “CRITICA DA RAZÃO PURA”, IMMANUEL KANT, PRETENDEU UMA SÍNTESE
CONCILIATÓRIA, AO SUGERIR A EXISTÊNCIA DO QUE CHAMOU DE “JUÍZOS SINTÉTICOS A PRIORI”,
QUE NÃO DEPENDERIAM DA EXPERIÊNCIA, E DOS QUAIS, O TEMPO E AS FORMAS GEOMÉTRICAS DA
GEOMETRIA EUCLIDIANA SERIAM EXEMPLOS; COM A DESCRIÇÃO, À PARTIR DE LOBACHEVSKI, EM
1829, DE GEOMETRIAS NÃO EUCLIDIANAS, E O ADVENTO DA TEORIA DA RELATIVIDADE, ADMITINDO
A OPÇÃO DE O ESPAÇO-TEMPO FÍSICO SER DESCRITO DE FORMA ATÉ MAIS REALISTA, POR ALGUMA
DELAS, E O TEMPO DEPENDER DA VELOCIDADE DO OBSERVADOR, SUA TEORIA NÃO MANTEVE A
SUSTENTAÇÃO, A PONTO DO MATEMÁTICO FRANCÊS H. POINCARÉ, NO FINAL DO SÉCULO XIX,
TER AFIRMADO QUE: “AS LEIS MAIS GERAIS DA NATUREZA, COMO AS LEIS GERAIS DO ESPAÇO E
DO TEMPO, NÃO SÃO DERIVÁVEIS DA EXPERIÊNCIA E NEM VERDADES LÓGICAS, MAS
CONVENÇÕES USADAS NA SISTEMATIZAÇÃO DOS DADOS EMPÍRICOS”.
- COM A PUBLICAÇÃO, EM 1859, POR C. DARWIN, DA OBRA “A ORIGEM DAS ESPÉCIES”, TEVE INÍCIO
O PONTO DE VISTA EVOLUCIONISTA, SEGUNDO O QUAL, “SOMOS PRODUTOS DE UMA EVOLUÇÃO,
QUE ENVOLVEU UMA CONTÍNUA INTERAÇÃO COM O AMBIENTE”.
“Muitas vezes ouvimos dizer que ‘um único homem não pode mais abranger um campo bastante
amplo’ e que ‘há demasiada especialização estreita’...Necessitamos de uma abordagem mais
simples, mais unificada dos problemas científicos, necessitamos de homens que pratiquem a ciência
e não de uma ciência particular, numa palavra, necessitamos de generalistas científicos” (Bode e
col. 1949). (...especialmente para integrar equipes multidisciplinares!).
9
>PROGRAMAÇÃO MATEMÁTICA: LINEAR, CONVEXA, INTEIRA, OTIMIZAÇÃO; .
>MÉTODOS COMPUTACIONAIS (OU NUMÉRICOS); .
>MÉTODOS ESTATÍSTICOS: PROBABILIDADES, PREDIÇÕES PROBABILÍSTICAS, ESTATÍSTICA, TEORIA
DA DECISÃO ESTATÍSTICA, ETC, E .
- APÓS A CONSTATAÇÃO DAS LIMITAÇÕES TANTO DAS EXPERIMENTAÇÕES (DESDE O NÍVEL DAS
PARTÍCULAS ELEMENTARES-PRINCÍPIO DA INCERTEZA-, ATÉ O DO ESPAÇO SIDERAL), QUANTO DAS
TEORIAS (CONSTATADAS INCLUSIVE PELOS TEOREMAS DE GÖDEL), APRESENTOU ADERÊNCIA
CRESCENTE, O ENTENDIMENTO OPERACIONAL PROVISÓRIO, DE QUE: NO ESTÁGIO ATUAL DE
INFORMAÇÃO E CONHECIMENTO, PODEMOS AFIRMAR, COM PRECISÃO CRESCENTE, ESTARMOS
LOCALMENTE (NO ESPAÇO E NO TEMPO), NUMA CONFIGURAÇÃO DE RELATIVA ESTABILIDADE
(INCLUSIVE AMBIENTAL, O QUE JUSTIFICARIA NÃO SÓ A EXISTÊNCIA DAS LEIS FÍSICAS, MAS
TAMBÉM DA NOSSA PRÓPRIA EXISTÊNCIA E EVOLUÇÃO), E QUE, AS TEORIAS DE QUE DISPOMOS
REPRESENTAM APROXIMAÇÕES ÚTEIS DESSA NATUREZA LOCAL. //
10
I- CRONOLOGIA ANTIGA
~ 3100 AC : É CRIADO NA SUMÉRIA, O CALENDÁRIO LUNAR DE DOZE MESES, E UM
3,1415926 E 3,1415927 ;
11
II- O SURGIMENTO DA CIÊNCIA GREGA
1-SEGUNDO HISTORIADORES, À PARTIR DO SÉCULO XII AC, COMO RESULTADO DAS INVASÕES
DÓRICAS AOS REINOS MICÊNICOS OU AQUEANOS, CUJA CIVILIZAÇÃO SE DESENVOLVEU EM
ESTREITA LIGAÇÃO COM A CIVILIZAÇÃO CRETENCE E EM COMUNICAÇÃO COM OS POVOS
ORIENTAIS, MUITOS AQUEUS FORAM FORÇADOS A IMIGRAR PARA AS ILHAS E COSTAS DA ÁSIA
MENOR, ONDE FORAM FUNDADAS CIDADES GREGAS COMO MILETO, ÉFESO, ETC., QUE SE
TRANSFORMARAM EM GRANDES CENTROS ECONÔMICOS E CULTURAIS, FACILITADO PELA
COMUNICAÇÃO COM OS OUTROS POVOS, DECORRENTE DO DESENVOLVIMENTO DA NAVEGAÇÃO E
DA COMERCIALIZAÇÃO DA LÃ E PELE DE CARNEIRO, AZEITE DE OLIVA E OUTRAS MERCADORIAS;
DURANTE O SÉCULO VII AC, O ADVENTO DA MOEDA FACILITANDO AS TROCAS, ACELEROU O
DESENVOLVIMENTO DA REGIÃO;
3-SEGUNDO UMA TRADIÇÃO, QUE REMONTA AOS PRÓPRIOS GREGOS, O PRIMEIRO FILÓSOFO
GREGO, TERIA SIDO TALES DE MILETO (DE FAMÍLIA FENÍCIA), QUE TERIA VIVIDO ENTRE O FINAL DO
SÉCULO VII E MEADOS DO SÉCULO VI AC. SEGUNDO PAUL TANNERY, TALES INTRODUZIU NA
GRÉCIA, NOÇÕES DA MATEMÁTICA ORIENTAL, QUE ELE DESENVOLVEU E APERFEIÇOOU, E DE MITOS
EGÍPCIO: PARA ELE, A AGUA SERIA A FONTE PRIMÁRIA DE TUDO O QUE EXISTE;
“TODAS AS COISAS PODEM SER REDUZIDAS À SUA ESSÊNCIA NUMÉRICA (ATRAVÉS DOS NÚMEROS
RACIONAIS, OU SEJA, QUOCIENTES DE NÚMEROS INTEIROS)”;
- FORAM ENCONTRADOS FRAGMENTOS ESCRITOS ATRIBUÍDOS A UM MEMBRO DA ESCOLA
12
PITAGÓRICA, RELATANDO EXPERIMENTOS QUE BUSCAVAM RELACIONAR OS SONS PRODUZIDOS
POR DISCOS METÁLICOS DE DIVERSOS DIÂMETROS E POR GARRAFAS CONTENDO MAIS OU MENOS
LÍQUIDO, A FRAÇÕES NUMÉRICAS, OU SEJA, A EXTENÇÃO PARA DUAS E TRÊS DIMENSÕES, DO QUE
“JÁ SE SABIA” PARA AS CORDAS: QUE “OS CHAMADOS SONS HARMÔNICOS CORRESPONDEM A
FRAÇÕES NUMÉRICAS DA CORDA”;
13
DISTINGUIR QUANTIDADES, O QUE NÃO REQUER ESPECÌFICAMENTE NÚMEROS, DE MODO QUE A
COMPOSIÇÃO DIGITAL DE SONS E IMAGENS, ORA CONQUISTADA, PODE NÃO SER A PRETENDIDA
POR PITÁGORAS.
11- SÉCULO V AC- PARMÊNIDES, DE ELÉIA: “OS HOMENS TERIAM DOIS CAMINHOS:
- O DA VERDADE, CONDUZIDO PELA RAZÃO (“O QUE EXISTE, OU SEJA, O SER, É!”; “O SER E O
PENSAR SÃO A MESMA COISA”; “O QUE É, É!, E NÃO PODE DEIXAR DE SER”; “O SER É ETERNO,
IMÓVEL, FINITO, IMUTÁVEL, PLENO, CONTÍNUO, HOMOGÊNEO E INDIVISÍVEL”); E
- O DA OPINIÃO, INDUZIDA À PARTIR DA OBSERVAÇÃO SENSÍVEL, QUE PERCEBE UM MUNDO DE
COISAS DIVERSAS, MUTÁVEIS, E QUE NÃO CONDUZ À VERDADE”,
AQUILES E A TARTARUGA:
14
13- SÉCULO V AC – ANAXÁGORAS, DE CLAZÔMENA: “NADA NASCE NEM MORRE, SENÃO QUE, A
PARTIR DAQUILO QUE EXISTE, CRIAM-SE COMBINAÇÕES E SEPARAÇÕES”; “TODAS AS COISAS
ESTAVAM JUNTAS, INFINITAS AO MESMO TEMPO EM NÚMERO E PEQUENEZ, E NENHUMA DELAS
PODIA SER DISTINGUIDA, PORQUE A PEQUENEZ ERA, TAMBÉM, INFINITA”.
14- COSTUMA SER ATRIBUÍDO A LEUCIPO (SÉCULO V AC) E DEMÓCRITO (SÉCULO IV AC),
-HIPÓCRATES, DE COS: QUE FUNDOU A ESCOLA DE COS, ONDE SE LANÇARAM AS BASES DO ESTUDO
OBJETIVO DA MEDICINA, DENTRE AS QUAIS, A DE QUE A MEDICINA SE CONSTITUÍA EM DISCIPLINA
DISTINTA DA FILOSOFIA, E QUE NÃO PODIA SER ESTUDADA À PARTIR DE PRESSUPOSTOS, MAS
APENAS TER COMO BASE A OBSERVAÇÃO CRITERIOSA DOS FENÔMENOS, OU SEJA, A
EXPERIÊNCIA, GERANDO A LINHA DE ENTENDIMENTO DOS “EMPÍRICOS”;
JÁ OUTROS DOIS DESTACADOS MÉDICOS ERASÍSTRATO E HERÓFILO, QUE DISSECAVAM CADÁVERES
15
À BUSCA DA DESCRIÇÃO DA ANATOMIA E FISIOLOGIA HUMANAS, ERAM TRATADOS PELOS SEUS
RIVAIS COMO “RACIONALISTAS”;
- PARA HIPÓCRATES, O CORAÇÃO TERIA O PAPEL DE CENTRO DAS SENSAÇÕES; ELE ENTENDIA QUE
A SAÚDE DEPENDIA DO EQUILÍBRIO ENTRE OS QUATRO HUMORES: BÍLIS PRETA, BÍLIS AMARELA,
SANGUE E FLEUGMA; QUANDO OCORRESSE UM DESIQUILÍBRIO, O PRÓPRIO CALOR VITAL (ENERGIA)
DO CORPO, PROVOCARIA A EXPULSÃO DO HUMOR EM EXCESSO ATRAVÉS DAS VIAS EXCRETORAS
(URINA, FEZES, VÔMITO, SUOR E EXPECTORAÇÕES E EXCREÇÕES BUCAIS OU NASAIS), DE MODO
QUE, QUANDO A EXPULSÃO NÃO OCORRIA PELAS VIAS NATURAIS, PODERIA OCORRER
COMPLICAÇÕES, JUSTIFICANDO UMA LINHA DE TRATAMENTO BASEADO EM ERVAS FERVIDAS, COM
PROPRIEDADES LAXATIVAS, DIURÉTICAS OU EMETIZANTES (VOMITÓRIOS);
17- 508 AC- UMA REVOLTA POPULAR LIDERADA POR CLÍSTENES, INSTAUROU, EM ATENAS, O
GOVERNO DE UMA “ASSEMBLÉIA DE CIDADÃOS” , QUE NÃO INCLUIA ESTRANGEIROS, ESCRAVOS E
MULHERES, QUE SE CONSTITUÍAM NA MAIORIA DA POPULAÇÃO, E QUE PASSOU A DECIDIR OS
DESTINOS DA POLIS (CIDADE-ESTADO);
19- SÉCULO V AC- O FILÓSOFO SÓCRATES, QUE ATRAVÉS DE PERGUNTAS A SEUS INTERLOCUTORES,
DESENVOLVEU O “QUESTIONAMENTO SOCRÁTICO” (QUE, SEGUNDO UMA INTERPRETAÇÃO, TERIA
SIDO INSPIRADO NO MÉTODO DO SOFISTA GÓRGIAS, QUE NEGAVA TODO TIPO DE VERDADE
OBJETIVA), CRITICAVA A FILOSOFIA E O CONHECIMENTO DA ÉPOCA, ESPECIALMENTE “OS DOGMAS
E CRENÇAS A PRIORI”, INCLUSIVE COM FRASES FAMOSAS COMO “TUDO O QUE SEI É QUE NADA
SEI”, E “O ARTÍFICE SABE MAIS DO SEU OFÍCIO DO QUE O FILÓSOFO DA SUA FILOSOFIA”, ACHAVA,
CONFORME MENCIONADO NUM DOS DIÁLOGOS DE PLATÃO, QUE HAVIAM DADO POUCA
IMPORTÂNCIA AO DESCOBRIMENTO DOS NÚMEROS IRRACIONAIS; SEU QUESTIONAMENTO TERIA
INSPIRADO A OBRA, OU REFORÇADO A MOTIVAÇÃO DE PLATÃO;
16
A MORTE DE SÓCRATES, POR JAQUES-LOUIS DAVID/1787
20- OBSERVAÇÃO: DE MODO ANÁLOGO AO QUE HAVIA OCORRIDO COM ANAXÁGORAS, QUE
APESAR DA AMIZADE DO GOVERNANTE PÉRICLES, FOI CONDENADO A DEIXAR ATENAS PELAS SUAS
IDÉIAS CONTRÁRIAS A RELIGIOSIDADE POPULAR E OFICIAL, SÓCRATES FOI CONDENADO A MORTE
EM 399 AC (SEGUNDO UMA INTERPRETAÇÃO, POR HAVER APENAS CRITICADO), O QUE TERIA
REFORÇADO O DESENCANTO DE PLATÃO (QUE TINHA PARENTES POLÍTICOS), PARA COM A POLÍTICA
E DEMOCRACIA DA ÉPOCA, E, AO MESMO TEMPO, MOTIVADO SUAS PESQUISAS (CIENTÍFICAS E
POLÍTICAS);
21- SÉCULO IV AC- ACADEMIA DE PLATÃO: EM 387 AC, PLATÃO FUNDOU EM ATENAS, “A
ACADEMIA”, ONDE HAVIA A INSCRIÇÃO “AQUI NÃO ENTRE QUEM NÃO SOUBER GEOMETRIA” COM
A INTENÇÃO PRECÍPUA DE PRODUZIR CONHECIMENTO SEGURO; A PRIMEIRA ETAPA DESSE
PROJETO, TERIA SIDO ENCAMINHADA POR PLATÃO E EUDOXO (A QUEM SE ATRIBUI A TEORIA DAS
PROPORÇÕES, E TAMBÉM A DIVISÃO DO ANO EM 365 DIAS + ¼ ) E RESULTADO NA CHAMADA
“GEOMETRIA EUCLIDIANA” , EM HOMENAGEM AO MATEMÁTICO FENÍCIO EUCLIDES (DE
ALEXANDRIA), NASCIDO EM TIRO (HOJE, LÍBANO), EM ~ 330 AC;
A PRINCIPAL CONTRIBUIÇÃO DA ACADEMIA TERIA SIDO: A APRESENTAÇÃO SINTÉTICA DA TEORIA,
PELO QUE É CONHECIDO COMO MÉTODO AXIOMÁTICO, NO QUAL: OS CONCEITOS PRIMITIVOS
SÃO ENUNCIADOS, AS PROPRIEDADES BÁSICAS QUE OS REGEM (PRINCÍPIOS) SÃO POSTULADAS, E A
PARTIR DELES, OS CONCEITOS SEGUINTES SÃO DEFINIDOS E AS OUTRAS PROPRIEDADES,
DEDUZIDAS. ASSIM, A PARTIR DOS CONCEITOS PRIMITIVOS DE PONTO, RETA, PLANO, ETC. ,
FORAM POSTULADAS AS SEGUINTES PROPRIEDADES (POSTULADOS DE CONSTRUÇÃO): PODE-SE:
1- DESENHAR UMA LINHA (SEGMENTO DE ) RETA DE UM PONTO A OUTRO;
17
2- PROLONGAR (INDEFINIDAMENTE) UM SEGMENTO DE RETA PARA ‘CONSTRUIR’ UMA RETA;
3- DESCREVER UM CÍRCULO, COM CENTRO E RAIOS DETERMINADOS;
4- TODOS OS ÃNGULOS RETOS SÃO IGUAIS; .
5- SE UMA RETA CORTA DUAS OUTRAS RETAS; DE MODO QUE A SOMA DOS ÂNGULOS INTERNOS DO
MESMO LADO, É MENOR DO QUE DOIS ÂNGULOS RETOS, AS DUAS RETAS, SE PROLONGADAS
INDEFINIDAMENTE, SE ENCONTRAM DO LADO EM QUE OS ÂNGULOS SÃO MENORES DO QUE DOIS
ÂNGULOS RETOS;
“...EIS O CAMINHO QUE SEGUI. COLOCO EM CADA CASO UM PRINCÍPIO, AQUELE QUE JULGO O
MAIS SÓLIDO, E TUDO O QUE PARECE ESTAR EM CONCORDÂNCIA COM ELE, ADMITO COMO
VERDADEIRO, ADMITINDO COMO FALSO, O QUE COM ELE NÃO CONCORDA” (PRINCÍPIO DA
DEDUÇÃO).
18
DETALHE DA PINTURA “A ESCOLA DE ATENAS”
República, de Platão
19
A INTERCEPTA)”, PRINCIPALMENTE POR IMPLICAR NUM “ESPAÇO INFINITO”, O QUE TAMBÉM NÃO É
INTUITIVO! , MAS À ÉPOCA, NÃO HAVIAM SIDO DESENVOLVIDAS TEORIAS ALTERNATIVAS, VISTO
QUE AS “GEOMETRIAS NÃO EUCLIDIANAS” FORAM FORMULADAS SOMENTE À PARTIR DO SÉCULO
XIX, E AS TEORIAS DE ESPAÇO FÍSICO (ESPAÇO-TEMPO) NÃO EUCLIDIANO, SOMENTE À PARTIR DO
INÍCIO DO SÉCULO XX, COM A TEORIA DA RELATIVIDADE;
20
ROTAÇÃO)”; A ÊLE É ATRIBUÍDA A FRASE: “DAI-ME UMA ALAVANCA E UM PONTO DE APOIO, E
MOVEREI A TERRA”
Fp é a força potente;
Fr é a força resistente;
BP é o braço potente; e
BR é o braço resistente.
DIZ-SE, NÃO SE SABE AO CERTO, SE VERÍDICA OU SUGESTIVAMENTE, QUE TAL PRINCÍPIO TERIA SIDO
INSPIRADO DURANTE UM BANHO DE BANHEIRA!
Se denotarmos por:
21
m a massa do corpo imerso,
V o volume do corpo imerso,
g a aceleração da gravidade,
I a força de impulsão.
FONTE: WIKIPÉDIA
22
Como mostrado por Arquimedes, a área do segmento parabólico na figura de cima é igual a
4/3 da do triângulo inscrito na figura de baixo.
Esta prova utiliza uma variação da série 1/4 + 1/16 + 1/64 + 1/256 + · · · cujo resultado é1/3.
23
-SÉC.II AC –HIPARCO, DE NICÉIA: APRIMOROU DADOS E MÉTODOS BABILÔNICOS COMO A
TRIGONOMETRIA, PARA A ABORDAGEM DE PROBLEMAS ASTRONÔMICOS; TAMBÉM SUSTENTOU
O HELIOCENTRISMO; À PARTIR DE UM ECLIPSE LUNAR, DESENVOLVEU UM MÉTODO GEOMÉTRICO
PARA CALCULAR A DISTÂNCIA ENTRE A TERRA E A LUA;
- O QUESTIONAMENTO TEÓRICO PROPOSTO POR RENÉ DESCARTES NO SÉCULO XVII, PARA VALIDAR
OS PRINCÍPIOS DO MÉTODO AXIOMÁTICO (QUE, POR SINAL, NÃO PREVALECEU);
- AS PERGUNTAS DO MÉTODO PSICANALÍTICO DE SIGMUND FREUD, NO INÍCIO DO SÉC.XX ;
- UMA NOVA NOÇÃO DE ALMA (PSIQUÊ), COMO SEDE DA CONSCIÊNCIA E DO CARÁTER;
- A IDÉIA DE QUE A CONSCIENTIZAÇÃO DAS PESSOAS (“CONHECE-TE A TI MESMO”) E DAS SUAS
ATIVIDADES PRÁTICAS, PODE COMPOR UMA MELHORIA SOCIAL (DE TODA ATENAS);
- A IRONIA, E ENFATIZADO AS DEMONSTRAÇÕES POR ABSURDO (AO TENTAR CONDUZIR SEUS
INTERLOCUTORES À CONTRADIÇÃO);
A) PROMOVEU A GEOMETRIA COMO A NOVA BASE DA MATEMÁTICA (AO INVÉS DOS NÚMEROS
RACIONAIS DE PITÁGORAS); TAL MUDANÇA TERIA SIDO ESSENCIALMENTE CONCEITUAL, VISTO QUE,
NOS PRÓPRIOS ELEMENTOS, PASSOU-SE A TRABALHAR COM NÚMEROS REAIS, NAS QUESTÕES
MÉTRICAS;
24
B) ELEIÇÃO DO MÉTODO DA GEOMETRIA (MÉTODO AXIOMÁTICO) COMO PADRÃO DE
APRESENTAÇÃO DE TEORIAS (QUE VEIO A TER SUCESSO NA APRESENTAÇÃO DE TEORIA DEDUTIVAS
COM LEIS QUANTITATIVAS- MATEMÁTICAS E FÍSICAS ASSENTADAS- APRESENTANDO RESTRIÇÕES
PARA OUTRAS TEORIAS;
25
- O PSICODRAMA (DE J.MORENO, À PARTIR DE 1921)
-----------------------xxx-------------------
26
III-SÉC. VIII/XII: CONTRIBUIÇÃO DOS CIENTISTAS ÁRABES/ISLÂMICOS:
27- A PARTIR DE 751, O PAPEL, OBTIDO DOS CHINESES, SUBSTITUIU O PERGAMINHO E O PAPIRO,
FACILITANDO A DIVULGAÇÃO EM GERAL;
- DURANTE A DINASTIA TURCA DOS ABÁCIDAS, NA SEGUNDA METADE DO SÉCULO VIII, MAMUM, FILHO DO
CALIFA HARUM AL-RACHID, FUNDOU EM BAGDAD, A “CASA DA SABEDORIA (BAIT AL-HIKMÁT)”, ONDE FORAM
TRADUZIDAS PARA O ÁRABE, AS OBRAS DE PLATÃO, EUCLIDES, GALENO, ARISTÓTELES, ARQUIMEDES E
PTOLOMEU, QUE FORAM AMPLIADAS PELOS ESTUDIOSOS SUBSEQUENTES; QUANDO, NO FINAL DO SÉCULO
XI, AS CIDADES ÁRABES DO SUL DA ESPANHA (QUE HAVIAM SIDO OCUPADAS À PARTIR DE 711) COMEÇARAM
A CAIR, AS OBRAS DE SUAS GRANDES BIBLIOTECAS FORAM TRADUZIDAS PARA O LATIM (ESPECIALMENTE EM
TOLEDO, ONDE FUNDARAM UM CENTRO DE TRADUÇÃO) E INVADIRAM A EUROPA CRISTÃ, CONTRIBUINDO
DECISIVAMENTE PARA O RENASCIMENTO EUROPEU; VÁRIOS DELES FORAM GENERALISTAS, TENDO SE
DEDICADO À QUÍMICA, ASTRONOMIA, MATEMÁTICA, FÍSICA, MEDICINA, ETC; DESTACAMOS AS
CONTRIBUIÇÕES DE:
28- JABIR (YABIR) IBN HAYYAN, ABU MUSSA DJAFAR AL KUHI (CONHECIDO NO OCIDENTE POR DJEBER, OU
GEBER): “SUMMA PERFECTIONIS”: TRATADO DE QUÍMICA MAIS ANTIGO CONHECIDO; SÍNTESE DE DIVERSAS
SUBSTÂNCIAS COMO OS ÁCIDOS SULFÚRICO, NÍTRICO, CLORÍDRICO, CÍTRICO, ACÉTICO, TARTÁRICO, ETC.;
PROCESSOS DE EVAPORAÇÃO, CONDENSAÇÃO, DESTILAÇÃO E CRISTALIZAÇÃO;
29- AL-KHWARIZMI, (FOI DIRETOR DA CASA DA SABEDORIA, E DEU O NOME A ALGARISMO): CONSIDERADO O
INICIADOR DA ÁLGEBRA, POR SUA OBRA “COMPÊNDIO DE CÁLCULO DE AL-JABR”, ADOTOU O SISTEMA HINDU
DE NUMERAÇÃO, COM DÍGITOS DE 0 A 9, E, POR ISSO CHAMADO SISTEMA DECIMAL DE NUMERAÇÃO, QUE
SIMPLIFICOU OS CÁLCULOS ARITMÉTICOS; .
- APÓS IBN AL-HAITHAM, CHAMANDO-SE A QUANTIDADE INCÓGNITA, POR UMA LETRA (POR EXEMPLO x); e
- ESCREVENDO-SE AS PROPRIEDADES CONHECIDAS, QUE DEVAM SER SATISFEITAS POR x (QUE SÃO
CHAMADAS EQUAÇÕES);
27
EXEMPLO, EM SÍMBOLOS, SE A EQUAÇÃO SE ESCREVER (1) a . x = b , com a,b números reais e a ≠ 0 ;
dividindo ambos os membros por a (ou, de modo equivalente, multiplicando por a -1 = 1/a , obtem-se a-1 ( a.x)
= a-1 .b , ou seja , x = a-1 . b
OBS: na abordagem da equação geral f(x) = y , no teorema da função inversa, a condição a≠0
se generaliza para f’(x0) ≠ 0 (derivada de f não nula no ponto x0 em consideração).
31- IBN AL-HAYTHAM (ALHAZEN): PARTINDO DA CONSTATAÇÃO DE QUE NOSSOS SENTIDOS NÃO SÃO
IMUNES AOS ERROS, E QUE TAMBÉM O CAMINHO DA RAZÃO NÃO ESTAVA CONDUZINDO À VERDADE, UMA
VEZ QUE MESMO OS MAIS SÁBIOS DIVERGIAM SOBRE DIVERSAS QUESTÕES (POR EXEMPLO OS PONTOS DE
VISTA OPOSTOS NA TEORIA DA VISÃO, DEFENDIDOS POR ARISTÓTELES E SEUS SEGUIDORES DE UM LADO, E
EUCLIDES E PTOLOMEU DO OUTRO), E QUE NÃO HAVIA DIVERGÊNCIAS SOBRE AS EQUAÇÕES MATEMÁTICAS
E PROVAS GEOMÉTRICAS, É CONSIDERADO O 1º CIENTISTA NA ACEPÇÃO MODERNA, AO SE UTILIZAR DA
“EXPERIMENTAÇÃO CONTROLADA” PARA VALIDAR AS PROPOSIÇÕES CIENTÍFICAS, ALÉM DE PREGAR A SUA
DESCRIÇÃO MATEMÁTICA, SEMPRE QUE POSSÍVEL, O QUÊ, REENUNCIADO POR GALILEU GALILEI, FICOU
CONHECIDO COMO “MÉTODO CIENTÍFICO”;
FOI O 1º A USAR LETRAS PARA DESIGNAR QUANTIDADES INCOGNITAS OU VARIÁVEIS NAS EXPRESSÕES
MATEMÁTICAS PARA ABORDAR PROBLEMAS GEOMÉTRICOS, E, POR ISSO, É CONSIDERADO O PRECURSOR DA
CHAMADA “GEOMETRIA ANALÍTICA” DESENVOVIDA NO SÉCOLO XVII POR RENÉ DESCARTES;
28
TAMBÉM LHE É ATRIBUÍDO O PRIMEIRO ENUNCIADO DO PRINCÍPIO DA INÉRCIA: “OS CORPOS CELESTES, NA
CONDIÇÃO DE OBJETOS REAIS, MOVEM-SE EM SUA COMBINAÇÃO, DE MODO CONTÍNUO E PERMANENTE,
VISTO NÃO HAVER NENHUM OBSTÁCULO, REPULSÃO OU IMPEDIMENTO”;
32- IBN SINA: DEU O NOME A “MEDICINA”, POR SUA OBRA “CANON” (AL-KANÚN), QUE FOI USADA COMO
TEXTO NAS ESCOLAS DE MEDICINA DA EUROPA POR MAIS DE QUATRO SÉCULOS, E CONSIDERADO O MAIOR
FILÓSOFO DA ÉPOCA, FORMULOU O CONCEITO DE QUANTIDADE DE MOVIMENTO LINEAR;
29
EXEMPLO: O conjunto dos pontos P = (x,y) para os quais a x + b y = c , representa, em
geral, uma reta no plano (x,y). As equações simultâneas
ax + by = e
cx + dy = f
Em 200 AC, os chineses já conheciam um método para resolver sistemas com duas equações
lineares com duas incógnitas. Hewitt menciona que as culturas antigas avançadas (Iraq, Egito, India
e Maia) também o conheciam.
Na tentativa de imitar o método de solução formal adotado para uma equação com uma
incógnita, pode-se reescrever o sistema na forma matricial (2) A X = B , onde A=
é a matriz dos coeficientes , X = , é a matriz incógnita, e B= é a matriz
dos segundos membros, com o produto de matrizes sendo definido a propósito.
A matriz inversa de A , quando existe (caso em que A é dita ser inversível), é denotada por A -1 , e
tem a propriedade de A-1 . A = Id = Matriz identidade = , tal que Id .X = X , de modo a, ao
multiplicar ambos os membros da equação (1) por A , obtermos A-1 . (A X) = A-1. B , ou seja X =
-1
A-1. B , e o sistema admite solução única. Ocorre que nem sempre o sistema admite solução (ou
seja, é possível), e, nesse caso, pode apresentar uma única solução (sistema determinado), ou
mesmo, infinitas soluções (sistema possível indeterminado). .
A teoria moderna dos SISTEMAS DE EQUAÇÕES LINEARES teve início em 1683, num trabalho do
japonês Seki Kowa , trazendo a idéia de determinante ,, e contou com as contribuições do alemão
Leibniz (1693), do escocês C. Maclaurin (~1729/48) , do suíço G. Cramer (1750), (regra de
Cramer), dos franceses E. Bézout, Laplace e Cauchy (1812), Gauss (~1800), Jordan (1887)
(método de eliminação de Gauss-Jordan), e do alemãoJacobi.
O Teorema de Rouché-Capelli, Fontené, Kronecker , Frobenius, fornece uma condição necessária
e suficiente para a existência, unicidade e multiplicidade de solução para o sistema linear
OBS: Sabemos que uma matriz quadrada 2x2, A é inversível, se e somente se, det A ≠ 0 ,sse, ‘uma
linha(ou coluna) não é múltipla da outra’ , condição que para sistemas se traduz como ‘uma
linha(ou coluna) não é combinação linear das outras’, e que na abordagem da equação geral
f(x)=y , x, y no teorema da função inversa, se generaliza para ‘determinante da
matriz das derivadas parciais (Jacobiana) no ponto em consideração ‘.
30
EVOLUÇÃO À PARTIR DO MÉTODO CIENTÍFICO:
IV- PRINCÍPIOS DA ÓTICA GEOMÉTRICA, LEI DE HOOKE E AS
LEIS DA MECÂNICA
COMO FOI MENCIONADO, NO SÉC. IX/X, IBN AL-HAYTHAM (ALHAZEN) LANÇOU AS BASES
DO MÉTODO CIENTÍFICO;
Onde:
31
A normal é a semi-reta perpendicular a superfície refletora.
Ângulo de incidência é o ângulo formado entre o feixe de luz que incide sobre o objeto e a
normal.
Ângulo de reflexão é o ângulo que a direção de um feixe de luz refletida faz com a normal.
Leis de Refração:
32
dado B, e que apresentam refração ou reflexão na superfície S de separação entre os meios, a
trajetória do raio de luz, é a que apresenta o TEMPO DE PERCURSO MÍNIMO entre A e B ;
39- ENTRE 1563 E 1601, O DINAMARQUÊS TYCHO BRAHE FEZ OBSERVAÇÕES DO SISTEMA
SOLAR, E OBTEVE OS DADOS MAIS PRECISOS DA ÉPOCA, E À BEIRA DA MORTE, TERIA DITO A
SEU ASSISTENTE, J. KEPLER: “NÃO ME DEIXE TER PARECIDO VIVER EM VÃO”;
41- 1609/19 – JOHANES KEPLER: BASEADO NOS DADOS DE TYCHO BRAHE, FORMULOU AS LEIS
DO MOVIMENTO DOS PLANETAS:
I. “CADA PLANETA DESCREVE UMA ÓRBITA ELÍPTICA (OVAL) EM TORNO DO SOL, QUE
OCUPA UM DOS FOCOS” :
ELIPSE
F1 , F2 : FOCOS
II. “O RAIO VETOR QUE VAI DO SOL AO PLANETA, VARRE ÁREAS IGUAIS EM INTERVALOS
DE TEMPO IGUAIS” :
33
A1=A2
a T
T2 ~ a3 , onde T é o período
e
a o semi-eixo maior
E POTÊNCIA EM PRODUTO,
E NO SÉCULO XVIII, EULER O RELACIONOU COM A EXPONENCIAL;
43- 1637, O FRANCÊS RENÉ DESCARTES, EM SEU “DISCURSO SOBRE O MÉTODO”, TENTOU DESCREVER
DE FORMA GERAL, O QUE SERIA O “MÉTODO AXIOMÁTICO”, MAS NÃO ATENDEU ÀS EXIGÊNCIAS
34
CIENTÍFICAS DE OBJETIVAÇÃO, AO SUGERIR QUE OS AXIOMAS DAS TEORIAS PUDESSEM SER
“VALIDADOS APENAS PELO QUESTIONAMENTO TEÓRICO, OU SEJA, QUANDO NÃO HOUVESSE MAIS
RESSALVAS, O QUE SERIA EQUIVALENTE A OBTER UNANIMIDADE” (O QUE, INCLUSIVE, O TERIA
CONDUZIDO A ALGUMAS CONCLUSÕES PRECIPITADAS E ENGANOS). Talvez tenha se baseado na
parte que segue, dos ensinamentos de Ibn al-Haitham: “The duty of the man who
investigates the writings of scientists, if learning the truth is his goal, is to make himself
an enemy of all that he reads, and,.. attack it from every side. He should also suspect
himself as he performs his critical examination of it, so that he may avoid falling into
either prejudice or leniency.” Descartes tentou uma descrição do mecanismo de ação
da força gravitacional, que não foi constatado pelas experiências. No seu método,
preconizou a divisão do todo em partes, para a análise em separado, para, em etapa
seguinte, promover a síntese (regra descartada pela abordagem sistêmica). Em sua
análise da circulação sanguínea, contribuiu com a afirmação que as artérias e veias
seriam tubos que carregavam nutrientes pelo corpo.
44- 1638 – GALILEU GALILEI: ENUNCIOU OS TRÊS PRINCÍPIOS DO SEU MÉTODO CIENTÍFICO:
1) OBSERVAÇÃO DOS FENÔMENOS, TAIS COMO ELES OCORREM, SEM QUE O CIENTISTA SE
DEIXE PERTURBAR POR PRECONCEITOS EXTRACIENTÍFICOS, DE NATUREZA RELIGIOSA OU
FILOSÓFICA;
2) PRINCÍPIO DA EXPERIMENTAÇÃO CONTROLADA: NENHUMA AFIRMAÇÃO SOBRE
FENÔMENOS NATURAIS, QUE SE PRETENDA CIENTÍFICA, PODE PRESCINDIR DA VERIFICAÇÃO
DE SUA LEGITIMIDADE, ATRAVÉS DA PRODUÇÃO DO FENÔMENO EM “DETERMINADAS
CIRCUNSTÂNCIAS” (DE OBJETIVIDADE);
3) AS LEIS NATURAIS DEVEM SER QUANTIFICADAS POR EXPRESSÕES MATEMÁTICAS;
45- 1671- O INGLÊS R. HOOKE ENUNCIOU A “ LEI DA ELASTICIDADE DOS CORPOS, OU LEI DE HOOKE”:
“ATÉ O CHAMADO LIMITE DE ELASTICIDADE, O ALONGAMENTO Δx SOFRIDO POR UMA MOLA
ELÁSTICA, É PROPORCIONAL À FORÇA F QUE O CAUSOU, OU SEJA,
Δx = C. F , ONDE C É UMA CONSTANTE DE PROPORCIONALIDADE, QUE DEPENDE DA MOLA, DE MODO
QUE “QUANTO MAIOR A FORÇA, MAIOR O ALONGAMENTO, E, CESSADA A FORÇA, CESSA TAMBÉM A
DEFORMAÇÃO, E A MOLA VOLTA À SUA POSIÇÃO INICIAL”; NESSE INTERVALO DIZEMOS QUE O
ALONGAMENTO VARIA LINEARMENTE COM A FORÇA, UMA VEZ QUE O SEU GRÁFICO É UM SEGMENTO
DE RETA; ENTRE O LIMITE DE ELASTICIDADE E O CHAMADO “LIMITE DE PLASTICIDADE” OU “PONTO DE
RUPTURA”, TEMOS AS DEFORMAÇÕES INELÁSTICAS, PARA AS QUAIS “CESSADA A FORÇA, A MOLA
AINDA PERMANECE COM ALGUMA DEFORMAÇÃO (DEFORMAÇÃO PERMANENTE), E À PARTIR DESSE
PONTO, A MENOR FORÇA PROVOCARIA A RUPTURA DA MOLA, GERANDO UMA DESCONTINUIDADE”;
35
Δx
F
F
Plot of applied force F vs. elongation X for a helical spring according to Hooke's law (red
line) and what the actual plot might look like (dashed line)
36
A INTENSIDADE DA FORÇA PODE SER REPRESENTADA PELO TAMANHO DA FLECHA.
- 3ª LEI- LEI DE AÇÃO E REAÇÃO: “A TODA AÇÃO (FORÇA) DE UM CORPO SOBRE OUTRO,
CORRESPONDE UMA REAÇÃO DO SEGUNDO SOBRE O PRIMEIRO, NA MESMA DIREÇÃO DA
AÇÃO, IGUAL EM INTENSIDADE E DE SENTIDO CONTRÁRIO”.
( FIGURA):
- LEI DA GRAVITAÇÃO UNIVERSAL: “MATÉRIA ATRAI MATÉRIA NA RAZÃO DIRETA DAS MASSAS
(ISTO É, QUANTO MAIORES AS MASSAS, MAIOR A ATRAÇÃO) E NA RAZÃO INVERSA DO
QUADRADO DA DISTÂNCIA ENTRE ELAS (OU SEJA, QUANTO MAIOR O QUADRADO DA
DISTÂNCIA, MENOR A ATRAÇÃO)”.
37
SOUBERMOS (OU SEJA, MANTER OS “PROBLEMAS ABERTOS”) PARA A ABORDAGEM DAS
GERAÇÕES FUTURAS”
38
, , ,
Funções trigonométricas
Linearidade
Regra do produto
Regra do quociente
ou,
Regra da Cadeia
, ou
39
Propriedade: Se F é uma primitiva de uma função contínua f, toda outra primitiva
de f tem a forma G(x) = F(x) + C, onde C é uma constante.
Integral Indefinida: Se f é uma função contínua, então a sua integral indefinida
existe, e é dada por
SE UM PONTO MATERIAL SE DESLOCA NO EIXO DOS x , A 2ª LEI PODE SER EXPRESSA POR:
m. x’’(t) = f(t, x(t), x’(t) ) , ONDE m É A MASSA , x(t) , x’(t) E x’’(t) , A POSIÇÃO,
VELOCIDADE E ACELERAÇÃO DO PONTO, E f , A RESULTANTE DAS FORÇAS EXTERNAS NA
DIREÇÃO DO EIXO DOS x , QUE, EM GERAL PODE DEPENDER DE t , x(t) E x’(t) ;
40
49- Em geral uma EQUAÇÃO DIFERENCIAL ORDINÁRIA DE 1ª ORDEM (maior
ordem de derivação que comparece na relação dada), fornece como SOLUÇÃO GERAL
(fórmula que expressa as soluções), uma família de curvas a um parâmetro ( que não se
intersectam, se valer a unicidade de solução passando por cada ponto considerado), e,
para eleger uma, entre as soluções, pede-se que satisfaça alguma condição adicional, por
exemplo, que passe por determinado ponto:
Por exemplo:
Por exemplo:
A solução geral da equação y’’ (t) = 0 , obtida através de duas integrações, representa a
família de curvas
y (t) = c1. t + c2 , que para cada c1 , representa uma família de retas paralelas com
coeficiente angular c1 , e, para cada c2 , o feixe de retas, no plano (t , y (t)) , que
intersectam o eixo y (t=0 na equação) no ponto de ordenada c2 ;
50- AO SUBSTITUIR O SÍMBOLO DE UMA FORÇA GENÉRICA POR UMA SÉRIE DE POTÊNCIAS,
NEWTON REVELA QUAL A AMPLITUDE QUE TINHA PARA O CONCEITO DE FUNÇÃO; EULER,
41
ADIANTANDO-SE ÀS PROVAS (QUE SÓ VIRIAM NO SÉCULO XIX), JÁ MANTEVE AS FUNÇÕES
β e γ NA FORMA INTEGRAL; FOURIER, EM 1811, ENUNCIOU QUE TODA FUNÇÃO PODERIA
SER EXPRESSA PELO SEU DESENVOLVIMENTO EM SÉRIES TRIGONOMÉTRICAS (SÉRIES DE
FOURIER), REVELANDO TAMBÉM QUAL A AMPLITUDE QUE TINHA PARA O CONCEITO DE
FUNÇÃO, CUJA DEFINIÇÃO SÓ VIRIA EM MEADOS DO SÉCULO; E CAUCHY, EM 1823, JÁ
CARACTERIZAVA FUNÇÃO CONTÍNUA EM TERMOS DE EPSLONS E DELTAS;
Forças de Atrito:
O atrito resulta
da interação
entre dois
corpos;
Num
automóvel, o
motor aciona as
rodas, que ao tentar girar, empurram o solo (fixo) para
traz, cuja reação, é uma força sobre as rodas, para
frente, que o faz andar; quanto maior a aderência dos
pneus no solo, menor o escorregamento, maior o
atrito, e maior é a força de reação propulsora sobre as
rodas;
Considere um corpo apoiado sobre uma superfície horizontal e rígida. Se o corpo receber a
ação de uma força f, devido às rugosidades surge a força de atrito.
As forças de atrito são contrarias ao movimento. Existem dois tipos de atrito estático e
cinético. Quando existe força atuando em um corpo mas ele não se move, o atrito é
denominado estático, quando existe força atuando num corpo e ele se move, o atrito é
denominado cinético.
42
Força de Atrito Estático
Se o corpo é puxado, porém não consegue escorregar na superfície, significa que ele recebeu a
ação de uma força de atrito que impede seu movimento. Essa força é denominada atrito
estático. Nesse caso:
F = FAE
A força de atrito estático tem um limite máximo, denominado tem um limite máximo,
denominado de força de atrito estático máximo.
FAEmax = μe . N
O coeficiente é um numero adimensional que depende das rugosidades da face do corpo que
está apoiada e da superfície de contato. Quanto mais áspero for o corpo ou a superfície maior
será o coeficiente.
A força de atrito estático pode variar de zero ate seu limite máximo, em função da intensidade
da força aplicada. Então o corpo so deslizará na superfície quando a força F vencer o atrito
estático.
Se o corpo está escorregando na superfície de apoio, significa que a força de atrito que age
nele é cinético ou dinâmico. A força de atrito cinético é dado por:
43
FAC = μc . N
O coeficiente é um numero adimensional que depende das rugosidades da face do corpo que
está apoiada e da superfície de contato. A força de atrito cinético é constante e não depende
da velocidade de escorregamento do corpo.
Na prática, verifica-se que é mais difícil tirar um corpo do repouso do que mantê-lo em
movimento:
μe ≥ μc
Força de Resistência do Ar
Se um corpo se movimenta através de um fluido (um gás, um líquido ou um vapor) surge uma força que
se opõe a esse movimento. Em se tratando do ar, essa força é chamada de força de resistência do ar.
Graças a essa resistência é que o paraquedas existe.
Quando um corpo está em movimento, ele sofre a ação de forças dissipativas, entre as quais podemos
citar o atrito e a resistência do ar.
Para o movimento de um corpo em contato com o ar (como a queda livre, o movimento de uma
motocicleta ou de um avião) com uma velocidade qualquer, a força da resistência do ar é dada por:
Fr = K . v 2
Onde k é uma constante que depende da forma do corpo e da área da secção transversal do corpo,
perpendicular à direção do movimento.
44
Em carros de fórmula 1, por exemplo, as formas aerodinâmicas diminuem o valor de K, o que ajuda a
diminuir a resistência do ar nesses veículos, fazendo com que ganhem mais velocidade.
Já nos paraquedas, por exemplo, sua aerodinâmica aumenta o valor de K, consequentemente a
resistência do ar aumenta.
Podemos dizer que o ar no paraquedas funciona como um vento forte, empurrando-o para cima, aliviando
a queda.
ou
45
onde:
, ou
Onde
é a velocidade da carga
O sentido da força magnética pode ser obtido pela regra da mão esquerda:
46
De modo equivalente, para o produto vetorial , é bastante usada a regra da mão
direita ou do saca-rolhas;
47
52-OCORRE QUE A 2ª LEI APRESENTOU OS SEGUINTES PROBLEMAS DE
APLICABILIDADE:
As equações de Euler-Lagrange
48
que representa um sistema de equações diferenciais parciais de 2ª ordem em t,
e L = T – V é a função lagrangeana , definida como a diferença entre a
energia cinética T e energia potencial V;
ou
, quando q=x
49
e as equações de Hamilton:
The HJE is also the only formulation of mechanics in which the motion of a particle can
be represented as a wave. In this sense, the HJE fulfilled a long-held goal of theoretical
physics (dating at least to Johann Bernoulli in the 18th century) of finding an analogy
between the propagation of light and the motion of a particle. The wave equation
followed by mechanical systems is similar to, but not identical with, Schrödinger's
equation, as described below; for this reason, the HJE is considered the "closest
approach" of classical mechanics to quantum mechanics. [1][2]
The dot (or inner) product notation between two lists of the same number of coordinates
is a shorthand for the sum of the products of corresponding components, e.g.
The dot product maps the two coordinate lists into one variable representing a single
numerical value.
50
where
is called Hamilton's principal function (also the action), qi are the N generalized
coordinates (i = 1,2...N) which define the configuration of the system, and t is time.
where the endpoints of the evolution are fixed and defined as and
. According to Hamilton's principle, the true evolution qtrue(t) is an
evolution for which the action is stationary. This principle results in the
equations of motion in Lagrangian mechanics.
The Hamilton's principal function S is related to the functional by fixing the initial
time t1 and endpoint q1 and allowing the upper limits t2 and the second endpoint q2 to
vary; these variables are the arguments of the function S. In other words, the action
function is the indefinite integral of the Lagrangian with respect to time.
The conjugate momenta correspond to the first derivatives of S with respect to the
generalized coordinates
51
As a solution to the Hamilton–Jacobi equation, the principal function contains N + 1
undetermined constants, the first N of them denoted as α1, α2 ... αN, and the last one
coming from the integration of .
The relationship between p and q then describes the orbit in phase space in terms of
these constants of motion. Furthermore, the quantities
are also constants of motion, and these equations can be inverted to find q as a function
of all the α and β constants and time.[4]
52
“NÃO É POSSÍVEL DETERMINAR PRATICAMENTE, AO MESMO TEMPO E COM PRECISÃO, A
POSIÇÃO E VELOCIDADE DE UMA PARTÍCULA ELEMENTAR: QUANTO MELHOR A PRECISÃO DE
UMA, MENOR A PRECISÃO DA OUTRA”;
A LIMITAÇÃO DO PROCESSO DE OBSERVAÇÃO É RELATIVAMENTE SIMPLES DE ENTENDER:
POR EXEMPLO, PARA OBSERVAR UMA MESA, FOTONS DE LUZ INCIDEM SOBRE A MESA, QUE
OS REFLETE AOS NOSSOS OLHOS; ESSE PROCESSO DE OBSERVAÇÃO PODE DESLOCAR
ALGUMAS PARTÍCULAS ELEMENTARES DA MESA, O QUE NÃO A ALTERA SIGNIFICATIVAMENTE;
PARA DETERMINAR A POSIÇÃO DE UMA PARTÍCULA ELEMENTAR COM PRECISÃO, DEVEMOS
INCIDIR SOBRE ELA UMA RADIAÇÃO COM COMPRIMENTO DE ONDA REDUZIDO (VISTO QUE,
QUANTO MENOR O COMPRIMENTO DE ONDA DA RADIAÇÃO INCIDENTE, MENOR A DISTÂNCIA
ENTRE DOIS PONTOS QUE PODEM SER DISTINGUIDOS E MAIOR A PRECISÃO NAS MEDIÇÕES DE
POSIÇÕES). OCORRE QUE, QUANTO MENOR O COMPRIMENTO DE ONDA DE UMA RADIAÇÃO,
MAIOR A SUA FREQUÊNCIA v E MAIOR A SUA ENERGIA E (EM VIRTUDE DA RELAÇÃO DE
PLANCK , ONDE h É A CONSTANTE DE PLANCK), E NA INCIDÊNCIA DESSA
RADIAÇÃO, HÁ A TRANSFERÊNCIA DE UMA MAIOR ENERGIA À PARTÍCULA OBSERVADA, O QUE
ALTERA SIGNIFICATIVAMENTE SUA VELOCIDADE v , E SEU MOMENTO LINEAR p= m.v (ONDE
m É SUA MASSA). COM RACIOCÍNIO ANÁLOGO, JUSTIFICA-SE QUE A MEDIÇÃO DA
VELOCIDADE DE UMA PARTÍCULA ELEMENTAR COM PRECISÃO, ALTERARIA
SIGNIFICATIVAMENTE A SUA POSIÇÃO;
DE ACORDO COM A MECÂNICA QUÂNTICA, ESSA INCERTEZA NÃO É UMA SIMPLES
LIMITAÇÃO DA NOSSA CAPACIDADE DE MEDIÇÃO, MAS FAZ PARTE DA “NATUREZA FÍSICA”;
DEPENDENDO DA FORMALIZAÇÃO, TRABALHA COM REGIÕES DE PROBABILIDADE DE CONTER
AS PARTÍCULAS ELEMENTARES;
MESMO QUE ESSE PRINCÍPIO PREVALEÇA AO NÍVEL DAS PARTÍCULAS ELEMENTARES, NADA
IMPEDE QUE HAJA ALGUMA PREDISIBILIDADE (PELO MENOS EM MÉDIA), AO NÍVEL DE
CONJUNTO DE PARTÍCULAS (MECÂNICA ESTATÍSTICA), COMO É O CASO DOS OBJETOS E
CORPOS DO NOSSO MUNDO REAL MACROSCÓPICO, OU DO SISTEMA SOLAR, QUE
APRESENTAM RELATIVA ESTABILIDADE;
- A MECÂNICA QUÂNTICA CONSEGUE JUSTIFICAR A ESTABILIDADE DOS ÁTOMOS, E, EM
PRINCÍPIO, PREVER A LIGAÇÃO ENTRE ÁTOMOS E AS REAÇÕES QUÍMICAS, REDUZINDO ASSIM,
A QUÍMICA À FÍSICA;
58- NA TENTATIVA DE AMPLIAR O ALCANCE DA TEORIA PARA OS CASOS NOS QUAIS NÃO SE
CONTAVA COM UMA FÓRMULA EXPLICITA PARA AS SOLUÇÕES DOS PROBLEMAS DA
MECÂNICA (P.EX. PROBLEMA DE n CORPOS SUJEITOS À ATRAÇÃO GRAVITACIONAL
MÚTUA, COMO PARA O SISTEMA SOLAR COM VARIOS PLANETAS), À PARTIR DA SEGUNDA
METADE DO SÉCULO XIX, TEVE INÍCIO A TEORIA ABSTRATA DAS EQUAÇÕES, COM
DEMONSTRAÇÕES DE EXISTÊNCIA E OUTRAS PROPRIEDADES DAS SOLUÇÕES, SEM A
NECESSIDADE DA EXPLICITAÇÃO DE SUAS FÓRMULAS ;
EXEMPLO: :
Picard–Lindelöf theorem, Picard's existence theorem or Cauchy–Lipschitz theorem is an
important theorem on existence and uniqueness of solutions to first-order equations with
given initial conditions, named after Charles Émile Picard, Ernst Lindelöf, Rudolf Lipschitz and
Augustin-Louis Cauchy:
53
Suppose is Lipschitz continuous in , |f(t,y1) - f(t,y2)| L | y1 -y2| , and continuous in
. Then, for some value , there exists a unique solution to the initial value
Peano existence theorem, Peano theorem (1890), named after Giuseppe Peano, is a
fundamental theorem which guarantees the existence of solutions to certain initial value
problems:
Observation: The Picard–Lindelöf theorem shows that the solution exists and that it is
unique. The Peano existence theorem shows only existence, not uniqueness, but it
assumes only that ƒ is continuous in y, instead of Lipschitz continuous. For example, the
right-hand side of the equation y′ = y1/3 with initial condition y(0) = 0 is continuous but
not Lipschitz continuous. Indeed, the solution of this equation is not unique; two
54
In the Observations, Peano wrote: “Consideróns en effet l’équation où
est continue et jamais nulle pour les valeurs de x dans l’intervalle (a,b). Alors, si est
il existe une seule function x de t, définie dans l’intervalle (a’, b’) qui satisfait à l’equation
donnée, et qui pour t=t0 a la valeur x0 ; comme cela résulte de líntégration de cette equation.
Et sur la derive de , on ná pas fait d’hypothèses.” (O resultado enunciado segue
According to V. I. Arnold, Ordinary Differential Equations, The MIT Press (1978), the idea behind the
theorem is that: beginning with an initial condition , . If the stationary
solution ,is reached after an infinite time the unicity of solution is guaranteed. However, if the
stationary solution is reached after a finite time, the unicity is violated
55
OBSERVAÇÃO: FORAM OBTIDAS ALGUMAS GENERALIZAÇÕES DOS RESULTADOS ACIMA, MAS
POR SEREM ELABORADAS E NÃO ALCANÇAREM SITUAÇÕES CONSIDERADAS RELEVANTES
(TENDO EM VISTA QUE “A LIMITAÇÃO DO ALCANCE JÁ ESTÁ INDICADA POR CONTRAEXEMPLOS”,
NÃO SÃO APRESENTADAS NOS TEXTOS OU RELEGADAS A “EXERCÍCIOS COM SUGESTÕES” (EX.
Theory of Ordinary Differential Equations, E.A. Coddington, N. Levinson);
-EM 1781, EM SUA “CRITICA DA RAZÃO PURA”, IMMANUEL KANT, PRETENDEU UMA SÍNTESE
CONCILIATÓRIA, AO SUGERIR A EXISTÊNCIA DO QUE CHAMOU DE “JUÍZOS SINTÉTICOS A
PRIORI”, QUE NÃO DEPENDERIAM DA EXPERIÊNCIA, E DOS QUAIS, O TEMPO E AS FORMAS
GEOMÉTRICAS DA GEOMETRIA EUCLIDIANA SERIAM EXEMPLOS; COM A DESCRIÇÃO, À PARTIR
DE LOBACHEVSKI EM 1829, DE GEOMETRIAS NÃO EUCLIDIANAS, E O ADVENTO DA TEORIA DA
RELATIVIDADE, ADMITINDO A OPÇÃO DE O ESPAÇO-TEMPO FÍSICO SER DESCRITO DE FORMA
ATÉ MAIS REALISTA, POR ALGUMA DELAS, E O TEMPO DEPENDER DA VELOCIDADE DO
OBSERVADOR, SUA TEORIA PERDEU A SUSTENTAÇÃO;
- COM A PUBLICAÇÃO, EM 1859, POR C. DARWIN, DA OBRA “A ORIGEM DAS ESPÉCIES”, TEVE
INÍCIO O PONTO DE VISTA EVOLUCIONISTA, SEGUNDO O QUAL, SOMOS PRODUTOS DE UMA
EVOLUÇÃO, QUE ENVOLVEU UMA CONTÍNUA INTEIRAÇÃO COM O AMBIENTE,
A PONTO DO MATEMÁTICO FRANCÊS H. POINCARÉ, NO FINAL DO SÉCULO XIX, TER
AFIRMADO QUE: “AS LEIS MAIS GERAIS DA NATUREZA, COMO AS LEIS GERAIS DO ESPAÇO E
DO TEMPO, NÃO SÃO DERAVÁVEIS DA EXPERIÊNCIA E NEM VERDADES LÓGICAS, MAS
CONVENÇÕES USADAS NA SISTEMATIZAÇÃO DOS DADOS EMPÍRICOS”.
56
Em uma esfera, a soma dos
ângulos dum triângulo não é igual
a 180°. Uma esfera não é um
espaço euclidiano, mas
localmente as leis da geometria
euclidiana são boas aproximções.
Num pequeno triângulo na face da
Terra, a soma dos ângulos é muito
próximo de 180°. Uma esfera
pode ser representada por uma
colecção de mapas de duas
dimensões, portanto uma esfera é
uma variedade.
61-EM 1915, O FÍSICO TEÓRICO ALEMÃO ALBERT EINSTEIN PUBLICOU A SUA TEORIA GERAL DA
RELATIVIDADE: UMA MASSA ACARRETA UMA “CURVATURA NO ESPAÇO A SUA VOLTA” (ASSIM COMO
UM PESO, COLOCADO SOBRE UMA MESA, TENDE A DEFLETÍ-LA, E O MOVIMENTO NATURAL (LIVRE DE
FORÇAS), DESCREVE AS CHAMADAS “GEODÉSICAS” (QUE SÃO AS CURVAS QUE GENERALIZAM AS RETAS
NOS ESPAÇOS CURVOS, OU QUE FORNECEM A DISTÂNCIA MÍNIMA ENTRE DOIS PONTOS PRÓXIMOS).
NO EXEMPLO DA SUPERFÍCIE ESFÉRICA, AS GEODÉSICAS SERIAM INTERSECÇÕES DE PLANOS QUE
PASSAM PELO CENTRO, COM A SUPERFÍCIE. ASSIM, OS PLANETAS DESCREVERIAM ÓRBITAS ELÍPTICAS
EM TORNO DO SOL, NÃO PORQUE SÃO “PUXADOS POR FORÇAS” EM SUA DIREÇÃO, MAS SIM, PORQUE
O SOL TORNARIA O ESPAÇO CURVO, E OS MOVIMENTOS NATURAIS NESSE ESPAÇO, SERIAM ELIPSES.
UTILIZANDO UMA ANALOGIA, SE COLOCÁSSEMOS UMA BOLA DE FUTEBOL DE SALÃO S (SOL) SOBRE
UMA MEMBRANA ELÁSTICA, ELA DEFORMARIA A MEMBRANA; SE ABANDONÁSSEMOS UMA BOLA DE
BILHAR T (TERRA) NAS PROXIMIDADES DE S, ELA SE MOVIMENTARIA NA DIREÇÃO DE S PELA
DEFORMAÇÃO DO ESPAÇO, E NÃO PELA ATRAÇÃO DE UMA FORÇA VEICULADA POR OUTRA PARTÍCULA
(UMA VEZ QUE SE COLOCÁSSEMOS AS DUAS BOLAS À MESMA DISTÂNCIA SOBRE UMA SUPERFÍCIE
PLANA RÍGIDA, ELAS NÃO SE MOVIMENTAM). POR ISSO, COSTUMA-SE DIZER QUE A RELATIVIDADE
GERAL “GEOMETRIZA A GRAVIDADE” :
57
Analogia, em 2 dimensões, para a curvatura do espaço-tempo causada por uma massa
G
Geodésicas num espaço curvo
58
MÉTODO CIENTÍFICO DE ALHAZEN/GALILEU, E O MÉTODO AXIOMÁTICO PARA A
APRESENTAÇÃO DAS TEORIAS DEDUTIVAS, E MOTIVOU AS PESQUISAS SUBSEQUENTES;
59
SONORAS, TRANSFERÊNCIA DE CALOR, DIFUSÃO TÉRMICA, SITUAÇÕES DE EQUILÍBRIO,
ESTADOS ESTACIONÁRIOS, E OUTROS, FORAM DESENVOLVIDAS AS “EQUAÇÕES
DIFERENCIAIS PARCIAIS”:
Exemplos:
Equação da onda (d’Alembert 1747, Euler 1748, D.Bernoulli 1753 e Lagrange 1759 ):
Onde é o Laplaciano e u = u (t , x , y , z )
60
A Equação de Poisson, também independente do tempo t , descreve o potencial
elétrico em eletrostática:
, com incógnita ψ
que aparece, por exemplo, na busca de soluções das equações das ondas por
variáveis separadas;
61
que no caso eletrostático em que f = 0 , é a energia potencial;
Uma interpretação análoga pode ser obtida para a temperatura de um corpo: dada a
distribuição de temperatura na fronteira de um corpo, determinar a distribuição da
temperatura do corpo, em equilíbrio térmico, entendendo que o calor flui no sentido oposto
ao do gradiente de temperatura.
Uma condição de existência pode ser encontrada integrando a equação em todo o domínio D e
aplicando a primeira identidade de Green:
1.
2.
3.
onde é o vetor unitário exterior normal.
62
- OBSERVAÇÃO: O TEOREMA DA DIVERGÊNCIA DE OSTROGRADSKY-GAUSS , PERMITIU O
EQUACIONAMENTO DO BALANÇO DE QUANTIDADES QUE FLUEM ATRAVÉS DE REGIÕES DO ESPAÇO, E
SUA EVENTUAL CONSERVAÇÃO (INCLUINDO ENERGIA, MASSA, MOMENTO LINEAR E ANGULAR, CARGA
ELÉTRICA, ETC.), ATRAVÉS DA QUANTIFICAÇÃO DOS FLUXOS DESSAS QUANTIDADES ATRAVÉS DAS SUAS
SUPERFÍCIES LIMÍTROFES, E DA CONSIDERAÇÃO DAS FONTES E SORVEDOUROS INTERNOS À REGIÃO;
TEOREMA DA DIVERGÊNCIA: O FLUXO de uma grandeza através de uma superfície S , pode ser
definido como a quantidade dessa grandeza que atravessa a superfície por unidade de tempo;
esse conceito é estendido para campos de vetores (campos de força, elétrico, magnético, etc.)
como a quantidade de linhas de campo que atravessa a superfície por unidade de tempo, que
também pode ser expresso como uma integral de superfície (soma dos fluxos elementares
através dos elementos de superfície); O divergente de um campo de vetores num ponto pode ser
interpretado como o fluxo através de superfícies próximas que o envolvem por unidade de
volume (pontualizado, ou seja, limite quando o volume tende ao ponto), e sua integral de
volume fornece o fluxo total nesse volume:
63
Forma geral das Leis de Conservação:
onde:
64
Dessa forma geral podemos obter as equações para a conservação da massa
(equação da continuidade), do momento linear, momento angular e energia:
Onde é a densidade de massa (ou carga), ou seja, massa (ou carga) por unidade de
volume, e v é a velocidade do fluido. No caso de um fluido incompressível, =
constante, não é uma função do tempo ou espaço, e a equação se reduz a:
Charge conservation. Vector calculus can be used to express the law in terms of
charge density ρ (in coulombs per cubic meter) and electric current density J (in
amperes per square meter):
65
The term on the left is the rate of change of the charge density ρ at a point. The term on
the right is the divergence of the current density J. The equation equates these two
factors, which says that the only way for the charge density at a point to change is for a
current of charge to flow into or out of the point. This statement is equivalent to a
conservation of four-current.
Charge conservation requires that the net current into a volume must necessarily equal the net
change in charge within the volume.
This yields
66
to the area, at right angles to that gradient, through which the heat flows, it is
possible to obtain the heat equation. Let u be the density of some quantity such as
heat, chemical concentration, etc. If V contained in U is any smooth
subregion, the rate of change of the total quantity within V equal the negative of
the flux through V:
67
- APÓS A ÁLGEBRA EXTERIOR DE GRASSMANN (QUE É UMA TEORIA DE TENSORES), O CÁLCULO
TENSORIAL FOI APRESENTADO EM 1892 POR G.RICCI-CURBASTRO, COM O TÍTULO “ABSOLUTE
DIFERENTIAL CALCULUS” E DIVULGADO POR RICCI e T.LEVI-CIVITA, ATRAVÉS DO TEXTO
CLÁSSICO “MÉTHODES DE CALCUL DIFFERENTIEL ABSOLUT ET LEURS APPLICATIONS” E
REPRESENTAM TRANSFORMAÇÕES LINEARES ENTRE ESCALARES, VETORES, OU OUTROS
TENSORES (POR EXEMPLO, O PRODUTO ESCALAR, O PRODUTO VETORIAL, AS
TRANSFORMAÇÕES LINEARES, ETC.), E DEVEM SER INDEPENDENTES DO PARTICULAR SISTEMA
DE COORDENADAS USADO PARA SUA REPRESENTAÇÃO; APRESENTA AMPLAS APLICAÇÕES EM
GEOMETRIA DIFERENCIAL (ESPECIALMENTE COM AS FORMAS DIFERENCIAS INTRODUZIDAS
POR E.CARTAN, E FORMAS QUADRÁTICAS COMO O TENSOR MÉTRICO, TENSOR DE
CURVATURA DE RIEMANN,ETC.), ELASTICIDADE (TENSOR DAS TENSÕES , DAS DEFORMAÇÕES,
ETC.), FLUIDODINÂMICA, RELATIVIDADE GERAL (TENSORES MÉTRICO, DE CURVATURA, DE
TENSÃO-ENERGIA), ETC.;
Definition. There are several approaches to defining tensors. Although seemingly different,
the approaches just describe the same geometric concept using different languages and at
different levels of abstraction.
As multidimensional arrays
Just as a scalar is described by a single number, and a vector with respect to a given basis is
described by an array of one dimension, any tensor with respect to a basis is described by a
multidimensional array. The numbers in the array are known as the scalar components of the
tensor or simply its components. They are denoted by indices giving their position in the array,
in subscript and superscript, after the symbolic name of the tensor. The total number of indices
required to uniquely select each component is equal to the dimension of the array, and is called
the order or the rank of the tensor.[Note 2] For example, the entries of an order 2 tensor T would
be denoted Tij, where i and j are indices running from 1 to the dimension of the related vector
space.[Note 3]
Just as the components of a vector change when we change the basis of the vector space, the
entries of a tensor also change under such a transformation. Each tensor comes equipped with
a transformation law that details how the components of the tensor respond to a change of
basis. The components of a vector can respond in two distinct ways to achange of
basis (see covariance and contravariance of vectors), where the new basis vectors are
expressed in terms of the old basis vectors as,
where Ri j is a matrix and in the second expression the summation sign was suppressed (a
notational convenience introduced by Einstein that will be used throughout this article). The
components, vi, of a regular (or column) vector, v, transform with the inverse of the matrix R,
where the hat denotes the components in the new basis. While the components, wi, of a
covector (or row vector), w transform with the matrix R itself,
The components of a tensor transform in a similar manner with a transformation matrix for each
index. If an index transforms like a vector with the inverse of the basis transformation, it is
called contravariant and is traditionally denoted with an upper index, while an index that
transforms with the basis transformation itself is called covariant and is denoted with a lower
index. The transformation law for an order-m tensor with n contravariant indices
and m−n covariant indices is thus given as,
68
Such a tensor is said to be of order or type (n,m−n).[Note 4]
This discussion motivates the
following formal definition:[9]
Definition. A tensor of type (n, m−n) is an assignment of a multidimensional array
As multilinear maps
A downside to the definition of a tensor using the multidimensional array approach is that it is
not apparent from the definition that the defined object is indeed basis independent, as is
expected from an intrinsically geometric object. Although it is possible to show that
transformation laws indeed ensure independence from the basis, sometimes a more intrinsic
definition is preferred. One approach is to define a tensor as a multilinear map. In that approach
a type (n,m) tensor T is defined as a map,
where V is a vector space and V* is the corresponding dual space of covectors, which is linear
in each of its arguments.
By applying a multilinear map T of type (n,m) to a basis {ej} for V and a canonical cobasis {εi}
for V*,
an n+m dimensional array of components can be obtained. A different choice of basis will yield
different components. But, because T is linear in all of its arguments, the components satisfy
the tensor transformation law used in the multilinear array definition. The multidimensional array
of components of T thus form a tensor according to that definition. Moreover, such an array can
be realised as the components of some multilinear map T. This motivates viewing multilinear
maps as the intrinsic objects underlying tensors.
Using tensor products
Main article: Tensor (intrinsic definition)
For some mathematical applications, a more abstract approach is sometimes useful. This can
be achieved by defining tensors in terms of elements of tensor products of vector spaces, which
in turn are defined through a universal property. A type (n,m) tensor is defined in this context as
an element of the tensor product of vector spaces, [12]
69
Using the properties of the tensor product, it can be shown that these components satisfy the
transformation law for a type (m,n) tensor. Moreover, the universal property of the tensor
product gives a 1-to-1 correspondence between tensors defined in this way and tensors defined
as multilinear maps.
In Cartesian coordinates, this reduces to the much simpler form (This can be seen a
special case of Lagrange's formula, see Vector triple product.)
For expressions of the vector Laplacian in other coordinate systems see Nabla in
cylindrical and spherical coordinates.
In Cartesian coordinates, this reduces to the much simpler form (This can be seen a
special case of Lagrange's formula, see Vector triple product.)
For expressions of the vector Laplacian in other coordinate systems see Nabla in
cylindrical and spherical coordinates.
Generalization
The Laplacian of any tensor field ("tensor" includes scalar and vector) is defined as
the divergence of the gradient of the tensor:
For the special case where is a scalar (a tensor of rank zero), the Laplacian takes on
the familiar form.
70
shown to be equivalent to the divergence of the expression shown below for the
gradient of a vector:
And, in the same manner, a dot product, which evaluates to a vector, of a vector by the
gradient of another vector (a tensor of 2nd rank) can be seen as a product of matrices:
An example of the usage of the vector Laplacian is the Navier-Stokes equations for a
Newtonian incompressible flow:
where the term with the vector Laplacian of the velocity field represents
the viscous stresses in the fluid.
Another example is the wave equation for the electric field that can be derived from the
Maxwell equations in the absence of charges and currents:
where
71
APÊNDICE III - RIEMANN/ LEBESGUE INTEGRATION
FUNDAMENTAL THEOREMS OF CALCULUS……………………………..[72,86]
72
subintervals", while in the Lebesgue integral, "one is in effect partitioning
the range of f".
As h approaches 0, it can be seen that the right hand side of this equation is
simply the derivative A′(x) of the area function A(x). The left-hand side of the
equation simply remains f(x). It can thus be shown, in an informal way, that f(x)
= A′(x)
73
The area shaded in red stripes can be estimated as h times f(x). Alternatively, if
the function A(x) were known, it could be computed as A(x + h) − A(x). These
two values are approximately equal, particularly for small h.
Formal statements
There are two parts to the theorem. Loosely put, the first part deals with the
derivative of an antiderivative, while the second part deals with the relationship
between antiderivatives and definite integrals.
First part
Then, F is continuous on [a, b], differentiable on the open interval (a, b), and
Corollary
74
Second part
The Second part is somewhat stronger than the Corollary because it does not
assume that f is continuous.
When an antiderivative F exists, then there are infinitely many antiderivatives
for f, obtained by adding to F an arbitrary constant. Also, by the first part of the
theorem, antiderivatives of f always exist when f is continuous.
The fundamental theorem can be generalized to curve and surface integrals in
higher dimensions and on manifolds. The most familiar extensions of the
Fundamental theorem of calculus in higher dimensions are the Divergence
theorem, the Gradient theorem, and the Stokes' theorem.
The gradient theorem
The gradient theorem states that if the vector field F is the gradient of some
scalar-valued function, then F is a conservative (i.e. path-independent) vector
field. This theorem has a powerful converse; namely, if F is a conservative
vector field, then F is the gradient of some scalar-valued function.[2] It is quite
straightforward to show that a vector field is path-independent if and only if the
integral of the vector field over every closed loop in its domain is zero. Thus the
converse can alternatively be stated as follows: If the integral of F over every
closed loop in the domain of F is zero, then F is the gradient of some scalar-
valued function.
Definitions. For some scalar field f : U ⊆ Rn → R, the line integral along a piecewise smooth
curve C ⊂ U is defined as
where r: [a, b] → C is an arbitrary bijective parametrization of the curve C such that r(a) and
r(b) give the endpoints of C and .
Line integral of a vector field. For a vector field F : U ⊆ Rn → Rn, the line
integral along a piecewise smooth curve C ⊂ U, in the direction of r, is
defined as
75
where · is the dot product and r: [a, b] → C is a bijective parametrization of the curve C
such that r(a) and r(b) give the endpoints of C. .
For example, the work done on an object moving through an electric or gravitational
field (W=F·s) have natural continuous analogs in terms of line integrals (W=∫CF·ds).
Nesse caso, é intuitivo que: o trabalho, num campo dissipativo de forças, dependa
do caminho.
The left side is a volume integral over the volume V, the right side is the surface
integral over the boundary of the volume V. The closed manifold ∂V is quite
generally the boundary of V oriented by outward-pointing normals, and n is the
outward pointing unit normal field of the boundary ∂V. (dS may be used as a
shorthand for n dS.) By the symbol within the two integrals it is stressed once
more that ∂V is a closed surface. In terms of the intuitive description above, the
left-hand side of the equation represents the total of the sources in the
volume V, and the right-hand side represents the total flow across the boundary
∂V.
76
Note that vectors may point into
or out of the
sphere.
Corollaries
77
KELVIN–STOKES THEOREM
This is the classical result first discovered by Lord Kelvin, who communicated
it to George Stokes in a letter dated July 2, 1850. Stokes set the theorem as a
question on the 1854 Smith's Prize exam, which led to the result bearing his
name. This special case is often just referred to as the Stokes' theorem in many
introductory university vector calculus courses and as used in physics and
engineering. It is also sometimes known as the curl theorem.
which relates the surface integral of the curl of a vector field over a surface Σ in
Euclidean three-space to the line integral of the vector field over its boundary.
The curve of the line integral, ∂Σ, must have positive orientation, meaning that
dr points counterclockwise when the surface normal, dΣ, points toward the
viewer, following the right-hand rule.
One consequence of the formula is that the field lines of a vector field with zero
curl cannot be closed contours.
Here d is the exterior derivative, which is defined using the manifold structure
only. The theorem is often used in situations where M is an embedded oriented
submanifold of some bigger manifold on which the form is defined.
-------------------------------------------------xxx-----------------------------------------------
78
THE LEBESGUE INTEGRAL (named after Henri Lebesgue who introduced
the integral in 1904).
O Projeto de Lebesgue
In Lebesgue’s words: “it is our purpose to associate with every bounded function which is
defined in a finite interval (a,b) – positive negative or equal to zero- a certain finite number
which we will call the integral of f on (a,b)
(3)
(5)
The significance, necessity, and corollaries of the first five conditions of this problem of
integration are more or less evident...”
In 1905, under the assumption of the axiom of choice, the Italian mathematician
Giuseppe Vitali gave the first example of a subset of real numbers (the Vitali set) that
is not Lebesgue measurable, and in 1970, the American mathematician Robert M.
Solovay constructed Solovay's model, which shows that it is consistent with standard
set theory, excluding uncountable choice, that all subsets of the reals are measurable.
79
The Lebesgue integral of a positive function is defined by using the Lebesgue
measure of a set. It uses a Lebesgue sum where is the value of
the function in subinterval , and is the Lebesgue measure of the set
of points for which values are approximately .
where
Note that both f+ and f− are non-negative measurable functions. Also note that
If
or sometimes
The Lebesgue integral does not distinguish between functions which differ only
on a set of μ-measure zero. To make this precise, functions f and g are said to
be equal almost everywhere (a.e.) if
80
To wit, the integral respects the equivalence relation of almost-everywhere
equality.
Monotonicity: If f ≤ g, then
81
Let I be an interval in the real line R. A function f: I → R is absolutely
continuous on I if for every positive number , there is a positive
number such that whenever a finite sequence of pairwise disjoint sub-intervals
(xk, yk) of I satisfies
then
Equivalent definitions
The length of a vector x = (x1, x2, …, xn) in the n-dimensional real vector
space Rn is usually given by the Euclidean norm:
The Euclidean distance between two points x and y is the length of the
straight line between the two points.
For a real number p ≥ 1, the p-norm or Lp-norm of x is defined by
The L∞-norm or maximum norm (or uniform norm) is the limit of the Lp-norms
for . It turns out that this limit is equivalent to the following definition:
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For all p ≥ 1, the p-norms and maximum norm as defined above indeed
satisfy the properties of a "length function" (or norm), which are that:
. > the length of the sum of two vectors is no larger than the sum of lengths of
the vectors (triangle inequality).( )
Abstractly speaking, this means that Rn together with the p-norm is
a Banach space. This Banach space is the Lp-space over Rn.
Illustrations of unit circles in different p-norms (every vector from the origin to
the unit circle has a length of one, the length being calculated with length-
formula of the corresponding p).
Relations between p-norms
It is intuitively clear that straight-line distances in Manhattan are generally
shorter than taxi distances. Formally, this means that the Euclidean norm of any
vector is bounded by its 1-norm:
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For the opposite direction, the following relation between the 1-norm and the 2-
norm is known:
This inequality depends on the dimension n of the underlying vector space and
follows directly from the Cauchy–Schwarz inequality
In general, for vectors in where p > r > 0:
Here, a complication arises, namely that the series on the right is not always
convergent, so for example, the sequence made up of only ones, (1, 1, 1, …),
will have an infinite p-norm (length) for every finite p ≥ 1. The space ℓp is then
defined as the set of all infinite sequences of real (or complex) numbers such
that the p-norm is finite.
One can check that as p increases, the set ℓp grows larger. For example, the
sequence
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diverges for p = 1 (the harmonic series), but is convergent for p > 1.
One also defines the ∞-norm using the supremum:
and the corresponding space ℓ∞ of all bounded sequences. It turns out that [2]
if the right-hand side is finite, or the left-hand side is infinite. Thus, we will
consider ℓp spaces for 1 ≤ p ≤ ∞.
The p-norm thus defined on ℓp is indeed a norm, and ℓp together with this norm
is a Banach space.
Lp spaces
Let 1 ≤ p < ∞ and (S, Σ, μ) be a measure space. Consider the set of
all measurable functions from S to C (or R) whose absolute value raised to
the p-th power has finite integral, or equivalently, that
The set of such functions forms a vector space, with the following natural
operations:
follows from the inequality |f + g|p ≤ 2p-1 (|f|p + |g|p). In fact, more is
true. Minkowski's inequalitysays the triangle inequality holds for || · ||p. Thus the
set of pth power integrable functions, together with the function || · ||p, is
a seminormed vector space, which is denoted by .
This can be made into a normed vector space in a standard way; one simply
takes the quotient space with respect to the kernel of || · ||p. Since for any
measurable function f, we have that ||f||p = 0 if and only if f = 0 almost
everywhere, the kernel of || · ||p does not depend upon p,
For p = ∞, the space L∞(S, μ) is defined as follows. We start with the set of all
measurable functions from S to C (or R) which are essentially bounded, i.e.
bounded up to a set of measure zero. Again two such functions are identified if
85
they are equal almost everywhere. Denote this set by L∞(S, μ). For f in L∞(S, μ),
its essential supremum serves as an appropriate norm:
As before, we have
When p = 2; like the ℓ2 space, the space L2 is the only Hilbert space of this
class. In the complex case, the inner product on L2 is defined by
------------------------------------------
Obs. A version of the fundamental theorem of calculus holds for the Gâteaux
derivative of F, provided F is assumed to be sufficiently continuously
differentiable (Appendix VI) . Specifically:
------------------------x-------------------------
86
APÊNDICE IV – EQUAÇÕES PROVENIENTES DA FÍSICA .......[87,107]
LEIS DA TERMODINÂMICA:
Primeira lei da termodinâmica: Lei da conservação da energia: “Num sistema isolado a energia
interna permanece constante”, ou " a variação da energia interna de um sistema é igual à
diferença entre o calor e o trabalho (que são formas de transferência de energia) trocados pelo
sistema com o meio exterior."
W > 0: volume do sistema aumenta; o sistema realiza trabalho sobre a vizinhança (cujo
volume diminui).
87
Segunda Lei da Termodinâmica:
Na física e engenharia, define-se eficiência como sendo a relação entre a energia fornecida a
um sistema (seja em termos de calor ou de trabalho) e a energia produzida pelo sistema
(normalmente na forma de trabalho).. A eficiência de um processo é definida como:
onde
ELETROMAGNETISMO:
- No início do SÉC. XIX, H.C.OERSTED obteve evidência empírica da relação entre fenômenos
elétricos e magnéticos;
88
- Em 1831, M.FARADAY descobriu a indução eletromagnética mostrando que era possível
transformar energia mecânica em energia elétrica, na primeira demonstração de um dínamo
que veio a ser a principal forma de obtenção de energia elétrica: fez um disco de cobre (ligado
a um galvanômetro) girar entre os polos de um imã em forma de ferradura e constatou que a
agulha do galvanômetro se movia com o girar do disco.
A Lei de indução eletromagnética de Faraday, enuncia que: a corrente elétrica induzida em
um circuito fechado por um campo magnético, é proporcional ao número de linhas de
campo que atravessa a área envolvida pelo circuito, na unidade de tempo;
A lei de Faraday-Lenz enuncia que a força eletromotriz que é induzida em um circuito
elétrico é igual à variação do fluxo magnético no circuito.
A lei de indução eletromagnética de Faraday-Neumann-Lenz fornece o princípio base do
funcionamento dos alternadores, dínamos e transformadores.
EQUAÇÕES DE MAXWELL:
1) Lei de Gauss: descreve a relação entre um campo elétrico e as cargas elétricas geradoras do
campo:
Onde:
campo elétrico
89
O fluxo elétrico (integral de superfície do campo elétrico) por meio da
ou
ou
90
, Rede de corrente elétrica passando através da
superfície S (incluindo correntes livres e ligadas)
Do Teorema de Stokes:
Se C ( uma curva fechada no R3 ) é a fronteira de uma superfície S então para cada campo vetorial F
= [ F1 , F2 , F3 ], temos que
F . dS = rot (F ) . dA
C S
Resulta a interpretação física do rotacional em um ponto P, como um vetor cuja componente normal
91
IGUAL A ou
ou
onde:
= aceleração da gravidade
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As seguintes convenções precisam ser satisfeitas para que a equação se aplique:
Equação de Bernoulli:
onde:
Uma segunda forma, mais geral, pode ser escrita para fluidos compressíveis:
onde:
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é a energia potencial gravitacional por unidade de massa, que vale apenas
no caso de um campo gravitacional uniforme,
Aviões: A asa de um avião é mais curva na parte de cima. Isto faz com que o ar passe mais
rápido na parte de cima do que na de baixo. De acordo com a equação de Bernoulli, a pressão
do ar em cima da asa será menor do que na parte de baixo, criando uma força de empuxo que
sustenta o avião no ar:
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Condensation visible over the upper surface of a wing caused by the fall in temperature
accompanying the fall in pressure, both due to acceleration of the air.
Fonte: Wikipedia
Equações de Navier-Stokes:
Forma Geral: com a notação do ítem 54, a forma das equações para a conservação do
momento é:
onde
onde:
a 2ª parcela representa o tensor das tensões;
os são a tensão normal;
tensão tangencial (tensão de cisalhamento);
p é a pressão estática, associada como a parte isotrópica do tensor de tensões sem
considerar se o fluido está ou não em equilíbrio; e
ρf outras forças externas, por unidade de volume, sobre o fluido (por exemplo: força
gravitacional ρg , forçacentrífuga, etc. ), uma vez que o gradiente de pressão já se
encontra explicitado;
Finalmente, temos:
Esta equação descreve a conservação do momento (que pode ser derivada a partir da 2ª
lei de Newton) para um fluido, mas se aplica a todo contínuo não relativista, e é
conhecida como Equação do momento de Cauchy;
Observação: não há prova de existência geral para o contexto tridimensional das Equações de
Navier-Stokes, ou existindo, se possuem singularidades. O “Clay Mathematics Institute”
considerou este, um dos sete problemas abertos mais importantes em matemática, e ofereceu
um prêmio de US$ 1.000.000 por uma solução ou contraexemplo;
Dependendo do caso particular, pode haver teoremas de existência, e/ou soluções
numéricas aproximadas;
Fonte: Navier-Stokes Equations – Wikipedia;
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MECÂNICA DO CONTÍNUO ( iniciada por A.L.Cauchy no Séc. XIX ):
Embora saibamos que a matéria é formada por átomos e moléculas, para muitos propósitos a
aproximação contínua, mais simples, se mostra útil e conveniente;
Elasticity
Describes materials that return to their rest shape
Solid mechanics after an applied stress.
The study of the physics of
Continuum continuous materials with a Plasticity
mechanics defined rest shape. Describes materials that Rheology
The study of the permanently deform after The study of materials
physics of a sufficient applied stress. with both solid and fluid
continuous characteristics.
materials Fluid mechanics Non-Newtonian fluids
The study of the physics of
continuous materials which
take the shape of their Newtonian fluids
container.
onde:
Desigualdade de Clausius–Duhem:
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A quantidade de interesse, no caso, é a entropia; supondo que há uma densidade de entropia
por unidade de massa ( ) interna à região , um fluxo e uma fonte de entropia, a 2ª Lei da
termodinâmica enuncia que: a taxa de aumento de na região , é maior ou igual à soma
das entropias acrescentadas à , por meio do fluxo de entropia, fontes internas, ou da
densidade de entropia devido ao fluxo de matéria ( que entre ou sai ) na região:
Sendo:
a densidade de matéria na região;
o fluxo de entropia na superfície; e
a fonte de entropia por unidade de massa;
O fluxo escalar de entropia pode ser relacionado com o fluxo vetorial na superfície, pela
relação:
Onde:
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MAGNETO FLUID DYNAMICS or MAGNETOHYDRODYNAMICS (MHD), is an academic
discipline which studies the dynamics of electrically conducting fluids. Examples of
such fluids include plasmas, liquid metals, and salt water or electrolytes. The word
magnetohydrodynamics (MHD) is derived from magneto- meaning magnetic field,
hydro- meaning liquid, and -dynamics meaning movement. The field of MHD was
initiated by Hannes Alfvén in 1942, for which he received the Nobel Prize in Physics in
1970: "At last some remarks are made about the transfer of momentum from the Sun to
the planets, which is fundamental to the theory (§11). The importance of the
magnetohydrodynamic waves in this respect are pointed out."
The ebbing salty water flowing past London's Waterloo Bridge interacts with the Earth's
magnetic field to produce a potential difference between the two river-banks. Michael
Faraday tried this experiment in 1832 but the current was too small to measure with the
equipment at the time, and the river bed contributed to short-circuit the signal.
However, by the same process, Dr. William Hyde Wollaston was able to measure the
voltage induced by the tide in the English Channel in 1851.
The fundamental concept behind MHD is that magnetic fields can induce currents in a
moving conductive fluid, which in turn creates forces on the fluid and also changes the
magnetic field itself. The set of equations which describe MHD are a combination of the
Navier-Stokes equations of fluid dynamics and Maxwell's equations of
electromagnetism. These differential equations have to be solved simultaneously, either
analytically or numerically.
MECÂNICA QUÂNTICA:
- Em 1900, o alemão M.E.Planck, das seus estudos sobre a emissão de radiação dos corpos
negros, sugeriu que a energia dos osciladores não é emitida em forma contínua, mas em
quantidades múltiplas inteiras de h. ν , onde h é a constante de Planck e ν a frequência (cor)
da radiação emitida, ou seja E = n.h. ν , onde n = 1, 2, 3, ....
- Em 1911, surgiu o modelo planetário de E.Rutherford que com base na experiência da “folha
de ouro”, sustentou que os átomos têm sua carga positiva concentrada em um pequeno
núcleo, com os elétrons orbitando o núcleo atômico, o que não estava de acordo com o
eletromagnetismo clássico, que prediz que os elétrons, ao orbitar os núcleos dos átomos,
como cargas elétricas em movimento, deveriam emitir energia, e quase imediatamente caírem
no núcleo, o que não se constata;
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4. Os elétrons podem saltar de um nível para outro mais externo, desde que absorvam uma
quantidade bem definida de energia (quantum de energia).
5. Ao voltar ao nível mais interno, o elétron emite um quantum de energia, na forma de luz
de cor bem definida ou outra radiação eletromagnética (fóton).
6. Cada órbita é denominada de estado estacionário e pode ser designada por letras K, L,
M, N, O, P, Q. As camadas podem apresentar:
K = 2 elétrons
L = 8 elétrons
M = 18 elétrons
N = 32 elétrons
O = 32 elétrons
P = 18 elétrons
Q = 2 elétrons
7. Cada nível de energia é caracterizado por um número quântico (n), que pode assumir
valores inteiros: 1, 2, 3, etc.
100
-Em 1924, o físico francês L.de Broglie conseguiu associar uma vibração às partículas
elementares, que lhes atribuiria características ondulatórias, conseguindo explicar a condição
de quantização dos átomos de Bohr e Sommerfeld (as ondas associadas aos elétrons eram
estacionárias no átomo, e as órbitas compreendiam um número inteiro de comprimentos de
onda), e obedeciam às equações de Einstein e de Planck: E = h. ν = m. c²
-Em 1925, o físico austríaco E.Schrödinger, sugeriu que a dinâmica da física das partículas
quânticas resultasse da chamada equação de onda ou equação de Schrödinger;
-Em 1926, Schrödinger apresentou uma demonstração de que os dois formalismos eram
equivalentes;
-Em 1926, o físico teórico inglês P.Dirac usou o formalismo dos colchetes de Poisson,
inspirando a unificação dos formalismos anteriores, através da teoria dos operadores;
p = | Ψ |² = Ψ* . Ψ
de modo que a probabilidade de encontrar a partícula num volume V , é dada pela integral
-Em 1928, P.Dirac desenvolveu uma teoria quântica relativística para o elétron, que permitiu
compreender o spin, identificado por W.Pauli ; também conseguiu prever a existência do
pósitron, antipartícula do elétron, identificada por C.Anderson, em 1932. No mesmo ano, o
físico inglês J. Chadwick identificou o Neutron.
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onde:
Ψ é a função de onda da mecânica quântica, que fornece uma descrição (máxima possível)
de um estado de um sistema, e, portanto, o caracteriza
i a unidade imaginária ,
que expressa a taxa de variação da função de onda em relação ao tempo, como função da
energia do estado;
Equação não-relativistica para uma única partícula que se move num campo elétrico (mas
não magnético):
onde:
ou, substituindo r = (x , y , z) ,
Equação geral:
102
hamiltoniano é proporcional à própria função Ψ , e a constante de proporcionalidade
(autovalor de correspondente ao autovetor Ψ ) é a energia E do estado Ψ ;
e a equação não-relativística, para uma única partícula que se move num campo elétrico
(mas não magnético):
The basic idea of the path integral formulation can be traced back to Norbert Wiener,
who introduced the Wiener integral for solving problems in diffusion and Brownian
motion. This idea was extended to the use of the Lagrangian in quantum mechanics by
P. A. M. Dirac in his 1933 paper. The complete method was developed in 1948 by
Richard Feynman. Some preliminaries were worked out earlier, in the course of his
doctoral thesis work with John Archibald Wheeler. The original motivation stemmed
from the desire to obtain a quantum-mechanical formulation for the Wheeler-Feynman
absorber theory using a Lagrangian (rather than a Hamiltonian) as a starting point.
The path integral also relates quantum and stochastic processes, and this provided the
basis for the grand synthesis of the 1970s which unified quantum field theory with
the statistical field theory of a fluctuating field near a second-order phase
transition. The Schrödinger equation is a diffusion equation with an imaginary
diffusion constant, and the path integral is an analytic continuation of a method for
summing up all possible random walks. For this reason path integrals were used in the
study of Brownian motion and diffusion a while before they were introduced in
quantum mechanics.
Recently, path integrals have been expanded from Brownian paths to Lévy flights. The
Lévy path integral formulation leads to fractional quantum mechanics and a fractional
Schrödinger equation.
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TEORIA DA RELATIVIDADE:
Esse enunciado pode derivar da constatação de que: estando um astronauta numa nave
espacial, não há experimento (ou observação) que possa realizar, dentro da nave e sem olhar
para fora, de modo a concluir sobre seu estado de movimento retilíneo e uniforme em relação
a outro sistema de referência, por exemplo, à terra, e pode ser enunciada como:
1.As leis que governam as mudanças de estado em quaisquer sistemas físicos tomam a
mesma forma em todo sistema inercial de coordenadas.
Obsservação: O fato da velocidade da luz no vácuo ser a maior velocidade que uma partícula,
onda ou corpo poder atingir é coerente com o princípio do tempo mínimo de Fermat;
" Eu estava sentado em uma cadeira no escritório de patentes, em Berna, quando de repente
ocorreu-me um pensamento: se uma pessoa cair livremente, ela não sentirá seu próprio peso.
Eu estava atônito. Este simples pensamento impressionou-me profundamente. Ele me impeliu
para uma teoria da gravitação." (Albert Einstein)
104
acelerado por uma força F adequada, que imponha ao mesmo uma aceleração de
módulo igual mas de sentido contrário ao da aceleração gerada no primeiro caso pelo
campo de gravidade.
" dentro do âmbito da mecânica clássica, onde a priori são definidos de formas distintas, os
conceitos de massa inercial e massa gravitacional são equivalentes. "
mecânica de Newton;
105
onde além disso é o tensor de curvatura de Ricci, é o escalar de curvatura de
Ricci e é a constante cosmológica.
The equation above was formulated by Einstein as part of his groundbreaking general theory of
relativity in 1915. The theory revolutionized how scientists understood gravity by describing the
force as a warping of the fabric of space and time.
"It is still amazing that one such mathematical equation can describe what space-time is all
about," said Space Telescope Science Institute astrophysicist Mario Livio, who nominated the
equation as his favorite. "All of Einstein's true genius is embodied in this equation."
"The right-hand side of this equation describes the energy contents of our universe (including
the 'dark energy' that propels the current cosmic acceleration)," Livio explained. "The left-hand
106
side describes the geometry of space-time. The equality reflects the fact that in Einstein's
general relativity, mass and energy determine the geometry, and concomitantly the curvature,
which is a manifestation of what we call gravity."
"It's a very elegant equation," said Kyle Cranmer, a physicist at New York University, adding
that the equation reveals the relationship between space-time and matter and energy. "This
equation tells you how they are related — how the presence of the sun warps space-time so
that the Earth moves around it in orbit, etc. It also tells you how the universe evolved since the
Big Bang and predicts that there should be black holes."
----------------- xxx-------------------
OVERVIEW. At present, matter and energy are best understood in terms of the
kinematics and interactions of elementary particles. To date, physics has reduced the
laws governing the behavior and interaction of all known forms of matter and energy to
a small set of fundamental laws and theories. A major goal of physics is to find the
"common ground" that would unite all of these theories into one integrated theory of
everything, of which all the other known laws would be special cases, and from which
the behavior of all matter and energy could be derived (at least in principle).
107
In the standard model of particle physics, there are six quarks — fundamental
particles that are the building-blocks of many others. Each quark also has an anti-
matter partner, an anti-quark. Pairings of quarks and anti-quarks form 'mesons',
such as the K+; three quarks form 'baryons', such as the proton. The picture builds
up further: a three-quark proton and a three-quark neutron together form a
deuteron; adding more protons and neutrons — more three-quark combinations —
builds up atomic nuclei. The discoveries of what seem to be a new meson2, 3 and a
new baryon4 don't easily fit the established picture. The new meson may in fact be
a 'molecule' of two mesons, and the baryon might be a 'pentaquark' state.
The Standard Model of particle physics is a theory concerning the electromagnetic, weak, and
strong nuclear interactions, which mediate the dynamics of the known subatomic particles.
Developed throughout the mid to late 20th century, The Standard Model is truly “a tapestry
woven by many hands”, sometimes driven forward by new experimental discoveries, sometimes
by theoretical advances. It was a collaborative effort in the largest sense, spanning continents
and decades. The current formulation was finalized in the mid 1970s upon experimental
confirmation of the existence of quarks. Since then, discoveries of the bottom quark (1977), the
top quark (1995), and the tau neutrino (2000) have given further credence to the Standard
Model. More recently, (2011–2012) the apparent detection of the Higgs boson completes the set
of predicted particles.
The Standard Model falls short of being a complete theory of fundamental interactions
because it does not incorporate the full theory of gravitation as described by general
relativity, or predict the accelerating expansion of the universe (as possibly described by
dark energy). The theory does not contain any viable dark matter particle that possesses
108
all of the required properties deduced from observational cosmology. It also does not
correctly account for neutrino oscillations (and their non-zero masses). Although the
Standard Model is believed to be theoretically self-consistent, it has several apparently
unnatural properties giving rise to puzzles like the strong CP problem and the hierarchy
problem.
Elementary Particles
Types Generations Antiparticle Colors Total
Quarks 2 3 Pair 3 36
Leptons 2 3 Pair None 12
Gluons 1 1 Own 8 8
W 1 1 Pair None 2
Z 1 1 Own None 1
Photon 1 1 Own None 1
Higgs 1 1 Own None 1
Total 61
Particle classification: Fermions and Bosons (Note that mesons are bosons and hadrons; and
baryons are hadrons and fermions).
109
110
Summary of interactions between particles described by the Standard Model.
111
O MODELO PADRÃO
112
DE MODO INFORMAL, OS FÉRMIONS (ELÉTRONS, PÓSITRONS, NEUTRINOS, QUARKS,
PRÓTONS, NEUTRONS, ETC.) SÃO AS PARTÍCULAS QUE CONSTITUEM A MATÉRIA , OS BÓSONS
(FÓTONS, que mediam a interação eletromagnética, BÓSONS W e Z, que mediam a
interação fraca, ambas diminuindo com a distância entre as partículas, os GLÚONS,
que mediam a interação forte, responsáveis pela coesão dos núcleos atômicos e que
aumenta com a distância entre as partículas) SÃO AS PARTÍCULAS QUE TRANSMITEM AS
FORÇAS, E O BÓSON DE HIGGS SERIA RESPONSÁVEL PELA MASSA INERCIAL DAS OUTRAS
PARTÍCULAS.
O bóson de Higgs é o único bóson na teoria que não é um bóson de calibre; tem um status
especial na teoria, o que foi assunto de algumas controvérsias. Grávitons, os bósons que
acredita-se mediar a interação gravitacional, não é explicado no modelo padrão.
Testes e predições
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Os fermions podem ser agrupados em três gerações, a primeira consiste do elétron,
quark para cima e para baixo e o neutrino elétron. Toda a matéria ordinária é
feita desta primeira geração de partículas; as gerações mais altas de partículas
decaem rapidamente para a primeira geração e somente podem ser gerados por um curto
tempo em experimentos de alta-energia. A razão para este arranjo em gerações é que os
quatro fermions em cada geração comportam-se sempre exatamente como seus
contrapontos na outra geração; a única diferença e suas massas. Por exemplo, o elétron e
o muon têm sempre meio spin e carga elétrica unitária, mas o muon e cerca de 200
vezes mais massivo.
Historical background. The first step towards the Standard Model was Sheldon Glashow's
discovery in 1960 of a way to combine the electromagnetic and weak interactions. In 1967
Steven Weinberg and Abdus Salam incorporated the Higgs mechanism into Glashow's
electroweak theory, giving it its modern form.
The Higgs mechanism is believed to give rise to the masses of all the elementary particles in the
Standard Model. This includes the masses of the W and Z bosons, and the masses of the
fermions, i.e. the quarks and leptons.
The Higgs particle is a massive scalar elementary particle theorized by Robert Brout, François
Englert, Peter Higgs, Gerald Guralnik, C. R. Hagen, and Tom Kibble in 1964 (see 1964 PRL
symmetry breaking papers) and is a key building block in the Standard Model. It has no
intrinsic spin, and for that reason is classified as a boson (like the gauge bosons, which have
integer spin).
The Higgs boson plays a unique role in the Standard Model, by explaining why the other
elementary particles, except the photon and gluon, are massive. In particular, the Higgs
boson would explain why the photon has no mass, while the W and Z bosons are very
heavy. Elementary particle masses, and the differences between electromagnetism (mediated by
the photon) and the weak force (mediated by the W and Z bosons), are critical to many aspects
of the structure of microscopic (and hence macroscopic) matter. In electroweak theory, the
Higgs boson generates the masses of the leptons (electron, muon, and tau) and quarks. As the
Higgs boson is massive, it must interact with itself.
After the neutral weak currents caused by Z boson boson exchange were discovered at CERN in
1973, the electroweak theory became widely accepted and Glashow, Salam, and Weinberg
shared the 1979 Nobel Prize in Physics for discovering it. The W and Z bosons were discovered
experimentally in 1981, and their masses were found to be as the Standard Model predicted.
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The theory of the strong interaction, to which many contributed, acquired its modern form
around 1973–74, when experiments confirmed that the hadrons were composed of fractionally
charged quarks.
Higgs boson. Because the Higgs boson is a very massive particle and also decays almost
immediately when created, only a very high energy particle accelerator can observe and record
it. Experiments to confirm and determine the nature of the Higgs boson using the Large Hadron
Collider (LHC) at CERN began in early 2010, and were performed at Fermilab's Tevatron until
its closure in late 2011. Mathematical consistency of the Standard Model requires that any
mechanism capable of generating the masses of elementary particles become visible at energies
above 1.4 TeV; therefore, the LHC (designed to collide two 7 to 8 TeV proton beams) was built
to answer the question of whether the Higgs boson actually exists.
On 4 July 2012, the two main experiments at the LHC (ATLAS and CMS) both reported
independently that they found a new particle with a mass of about 125 GeV/c2 (about 133
proton masses, on the order of 10 −25 kg), which is "consistent with the Higgs boson." Although
it has several properties similar to the predicted "simplest" Higgs, they acknowledged that
further work would be needed to conclude that it is indeed the Higgs boson, and exactly which
version of the Standard Model Higgs it best supported if confirmed.
115
An example of simulated data modeled for the CMS particle detector on the Large Hadron
Collider (LHC) at CERN. Here, following a collision of two protons, a Higgs boson is produced
which decays into two jets of hadrons and two electrons. The lines represent the possible
paths of particles produced by the proton-proton collision in the detector while the energy
these particles deposit is shown in blue.
SCATTERING is a general physical process where some forms of radiation, such as light,
sound, or moving particles, are forced to deviate from a straight trajectory by one or
more localized non-uniformities in the medium through which they pass. In
conventional use, this also includes deviation of reflected radiation from the angle
predicted by the law of reflection. Reflections that undergo scattering are often called
diffuse reflections and unscattered reflections are called specular (mirror-like)
reflections.
Some areas where scattering and scattering theory are significant include radar sensing,
medical ultrasound, semiconductor wafer inspection, polymerization process
monitoring, acoustic tiling, free-space communications, and computer-generated
imagery.
Scattering theory is a framework for studying and understanding the scattering of waves and
particles. Prosaically, wave scattering corresponds to the collision and scattering of a wave
with some material object, for instance sunlight scattered by rain drops to form a rainbow.
Scattering also includes the interaction of billiard balls on a table, the Rutherford scattering (or
angle change) of alpha particles by gold nuclei, the Bragg scattering (or diffraction) of electrons
and X-rays by a cluster of atoms, and the inelastic scattering of a fission fragment as it
traverses a thin foil. More precisely, scattering consists of the study of how solutions of
partial differential equations, propagating freely "in the distant past", come together and
interact with one another or with a boundary condition, and then propagate away "to the
distant future".
Electromagnetic waves are one of the best known and most commonly encountered forms of
radiation that undergo scattering. Scattering of light and radio waves (especially in radar) is
particularly important. Several different aspects of electromagnetic scattering are distinct
enough to have conventional names. Major forms of elastic light scattering (involving
negligible energy transfer) are Rayleigh scattering and Mie scattering. Inelastic scattering
includes Brillouin scattering, Raman scattering, inelastic X-ray scattering and Compton
scattering.
116
Rayleigh and Mie Scattering
For modeling of scattering in cases where the Rayleigh and Mie models do not apply such as
irregularly shaped particles, there are many numerical methods that can be used. The most
common are finite-element methods which solve Maxwell's equations to find the distribution
of the scattered electromagnetic field. Sophisticated software packages exist which allow the
user to specify the refractive index or indices of the scattering feature in space, creating a 2- or
sometimes 3-dimensional model of the structure. For relatively large and complex structures,
these models usually require substantial execution times on a computer….
117
Standard model
Another of physics' reigning theories, the standard model describes the collection of
fundamental particles currently thought to make up our universe.
The theory can be encapsulated in a main equation called the standard model
Lagrangian (named after the 18th-century French mathematician and astronomer Joseph
Louis Lagrange), which was chosen by theoretical physicist Lance Dixon of the SLAC
National Accelerator Laboratory in California as his favorite formula.
"It has successfully described all elementary particles and forces that we've observed in
the laboratory to date — except gravity," Dixon told LiveScience. "That includes, of
course, the recently discovered Higgs (like) boson, phi in the formula. It is fully self-
consistent with quantum mechanics and special relativity."
The dynamics of the quantum state and the fundamental fields are determined by the
Lagrangian density (usually for short just called the Lagrangian). This plays a role similar to
that of the Schrödinger equation in non-relativistic quantum mechanics, but a Lagrangian is
not an equation – rather, it is a polynomial function of the fields and their derivatives. While it
would be possible to derive a system of differential equations governing the fields from the
Langrangian, it is more common to use other techniques to compute with quantum field
theories.
The standard model theory has not yet, however, been united with general
relativity, which is why it cannot describe gravity.
118
vibrantes, o que exige um aumento das dimensões do Espaço-Tempo); SITE OFICIAL DA
TEORIA DAS CORDAS: www.superstringtheory.com,
GRAVIDADE QUÂNTICA EUCLIDIANA, elaborada pelo físico britânico Stephen Hawking: numa
abordagem que postula que o espaço-tempo emerge de uma grande média quântica de todas
as possíveis formas;
Carlo Rovelli, Quantum Gravity, Cambridge University Press-2008- Lee Smolin, A Vida do
Cosmos, Editora Unisinos, São Leopoldo, 2004; Martin Bojowald- Living Reviews in Relativity,
vol.11,nº4 (July 2, 2008)- http:// www.livinggreviews.org/Articles/lrr-2008-4
119
g) OS MATEMÁTICOS BUSCARAM GENERALIZAR AO MÁXIMO, OS MÉTODOS UTILIZADOS
NA ABORDAGEM DESSAS EQUAÇÕES, SENDO OS MAIS UTILIZADOS, O DAS SÉRIES
(TRIGONOMÉTRICAS) DE FOURIER E O DA TRANSFORMADA DE FOURIER, COM OS PASSOS DE
TRANSFORMAR A EQUAÇÃO DIFERENCIAL NUMA EQUAÇÃO ALGÉBRICA, RESOLVÊ-LA, E, EM
SEGUIDA, USANDO A TRANSFORMAÇÃO INVERSA, OBTER UMA (FUNÇÃO OU DISTRIBUIÇÃO)
“CANDIDATA A SOLUÇÃO”;
120
= Id ( u) = u , onde Id representa o operador identidade , para, aplicando F -1 a
ambos os membros de (1) , obter u = F-1 (f).
O TEOREMA DA FUNÇÃO INVERSA fornece condições suficientes para que uma função
diferenciável seja inversível numa vizinhança de pontos do seu domínio. Há versões para o IR n ,
variedades, espaços de Banach , variedades de Banach e até mais gerais. In November 2010,
Ivar Ekeland presented “An Inverse Function Theorem in Frechet Spaces
ABSTRACT I present an inverse function theorem for differentiable maps between Frechet
spaces which contains the classical theorem of Nash and Moser as a particular case. In
contrast to the latter, the proof does not rely on the Newton iteration procedure, but on
Lebesgue's dominated convergence theorem and Ekeland's variational principle. As a
consequence, the assumptions are substantially weakened: the map F to be inverted is not
required to be C2, or even C1, or even Frechet-differentiable. Comment: to appear, Annales de
l'Institut Henri Poincare, Analyse Non Lineaire.”
À PARTIR DOS TEOREMAS ACIMA, PODE-SE DERIVAR TEOREMAS DAS FUNÇÕES IMPLÍCITAS
CORRESPONDENTES.
121
are encountered include radiative energy transfer and the oscillation of a string, membrane, or
axle. Oscillation problems may also be solved as differential equations.
Overview
The most basic type of integral equation is a Fredholm equation of the first type:
If the unknown function occurs both inside and outside of the integral, it is known as a
Fredholm equation of the second type:
The parameter λ is an unknown factor, which plays the same role as the eigenvalue in
linear algebra.
In all of the above, if the known function f is identically zero, it is called a homogeneous
integral equation. If f is nonzero, it is called an inhomogeneous integral equation.
Both Fredholm and Volterra equations are linear integral equations, due to the linear
behaviour of φ(x) under the integral. A nonlinear Volterra integral equation has the
general form:
122
OBSERVAÇÃO: USANDO A ADITIVIDADE DA INTEGRAL EM RELAÇÃO AOS LIMITES DE
INTEGRAÇÃO (EXTREMOS DO INTERVALO), É SIMPLES CONSTATAR QUE, SOB HIPÓTESES
SUFICIENTES, “SOLUÇÕES LOCAIS (EM SUBINTERVALOS) PODEM SER LIGADAS PARA SE
OBTER UMA SOLUÇÃO EM TODO O INTERVALO” :
g(x) = f (x) + ₁ (t )) dt
(7) obtemos 1 (a2) = ₂ (a2) , ou seja as funções têm o mesmo valor no ponto comum
aos intervalos.
EM TERMOS GERAIS, PODEMOS DIZER QUE OBTEVE-SE SEIS TIPOS BÁSICOS DE TEOREMAS EXISTÊNCIA,
AS VEZES UTILIZADOS EM COMBINAÇÃO, QUE ENVOLVEM HIPÓTESES MAIS OPERACIONAIS ou MAIS
SIMPLES DE SEREM VERIFICADAS, OU SUGERIDAS PELAS APLICAÇÕES:
II) TEOREMAS DE PONTO FIXO, NOS QUAIS, UTILIZANDO-SE O OPERADOR k DEFINIDO POR
123
k (u) = u + F (u) - f , REESCREVE-SE A EQUAÇÃO (1) , OBTENDO-SE:
(2) k(u) = u , E BUSCA-SE A SOLUÇÃO DA NOVA EQUAÇÃO, COMO UM PONTO FIXO DO OPERADOR k !
IV) TEOREMAS DO TIPO VARIACIONAIS, QUE GARANTEM A EXISTÊNCIA DE PONTOS EXTREMANTES (DE
MÁXIMO OU MÍNIMO) DE FUNÇÕES E OPERADORES (P. EX., BUSCANDO EQULÍBRIO EM
CONFIGURAÇÕES DE ENERGIA MÍNIMA, TRAJETÓRIAS DE TEMPO OU AÇÃO MÍNIMOS, ETC. P.EX.,
QUANDO F É UMA DERIVADA, F=A’, AS SOLUÇÕES DA EQUAÇÃO HOMOGÊNEA, F(u)=0, SÃO OS PONTOS
CRÍTICOS DE A);
j) CONVÉM NOTAR QUE APESAR DESSE ELABORADO APARATO TEÓRICO, MESMO QUANDO
UTILIZADO, A ABORDAGEM E ANÁLISE DE EQUAÇÕES ESPECÍFICAS, EM GERAL FORNECE
RESULTDOS ADICIONAIS;
EXEMPLOS:
t
(5) y(t) = y0 + f (s, y(s)) ds , que pode ser escrita como em
0
124
Ocorre que para a aplicação do teorema do ponto fixo de Banach, que enuncia que:
“todo operador num espaço de Banach (espaço vetorial normado e completo), para
o qual existe uma potência que é uma contração, admite um único ponto fixo nesse
espaço”, se faz necessária alguma hipótese adicional para f , por exemplo, a condição
de Lipschitz na 2ª variável
125
O resultado não é válido
para regiões “com
buracos” ou não limitadas,
como mostram os
exemplos ao lado;
126
APÊNDICE V - THEORY OF LINEAR PARTIAL DIFFERENTIAL EQUATIONS [127, 164]
we (formally) apply the Fourier transform on both sides and obtain the algebraic
equation
If the symbol P(ξ) is never zero when ξ ∈ Rn, then it is possible to divide by P(ξ):
The last assumption can be weakened by using the theory of distributions. The first two
assumptions can be weakened as follows.
This is similar to formula (1), except that 1/P(ξ) is not a polynomial function, but a
function of a more general kind.
127
naturally found in Sobolev spaces, rather than in spaces of continuous functions and
with the derivatives understood in the classical sense.
Motivation: There are many criteria for smoothness of mathematical functions. The
most basic criterion may be that of continuity. A stronger notion of smoothness is that
of differentiability (because functions that are differentiable are also continuous) and a
yet stronger notion of smoothness is that the derivative also be continuous (these
functions are said to be of class C1 — see smooth function). Differentiable functions are
important in many areas, and in particular for differential equations. On the other hand,
quantities or properties of the underlying model of the differential equation are usually
expressed in terms of integral norms, rather than the uniform norm. A typical example
is measuring the energy of a temperature or velocity distribution by an L2-norm. It is
therefore important to develop a tool for differentiating Lebesgue functions.
The integration by parts formula yields that for every u ∈ Ck(Ω), where k is a natural
number and for all infinitely differentiable functions with compact support φ ∈ Cc∞(Ω),
where α a multi-index of order |α| = k and Ω is an open subset in ℝn. Here, the notation
is used.
The left-hand side of this equation still makes sense if we only assume u to be locally
integrable. If there exists a locally integrable function v, such that
we call v the weak α-th partial derivative of u. If there exists a weak α-th partial
derivative of u, then it is uniquely defined almost everywhere. On the other hand, if u ∈
Ck(Ω), then the classical and the weak derivative coincide. Thus, if v is a weak α-th
partial derivative of u, we may denote it by Dαu := v.
The Sobolev spaces Wk,p(Ω) combine the concepts of weak differentiability and
Lebesgue norms.
Definition:
The Sobolev space Wk,p(Ω) is defined to be the set of all functions u ∈ Lp(Ω) such that
for every multi-index α with |α| ≤ k, the weak partial derivative belongs to Lp(Ω),
i.e.
128
Here, Ω is an open set in ℝn and 1 ≤ p ≤ +∞. The natural number k is called the order of
the Sobolev space Wk,p(Ω).
There are several choices for a norm for Wk,p(Ω). The following two are common and
are equivalent in the sense of equivalence of norms:
and
With respect to either of these norms, Wk,p(Ω) is a Banach space. For finite p, Wk,p(Ω)
is also a separable space. It is conventional to denote Wk,2(Ω) by Hk(Ω) for it is a
Hilbert space with the norm .
Generalized functions were introduced by Sergei Sobolev in 1935. They were re-
introduced in the late 1940s by Laurent Schwartz, who developed a comprehensive
theory of distributions.
A typical test function, the bump
Basic idea
function Ψ(x). It is smooth
(infinitely differentiable) and has
compact support (is zero outside
an interval, in this case the
interval [-1, 1]).
Distributions are a class of linear functionals that map a set of test functions
(conventional and well-behaved functions) onto the set of real numbers. In the simplest
case, the set of test functions considered is D(R), which is the set of functions from R to
R having two properties:
129
The function is smooth (infinitely differentiable);
The function has compact support (is identically zero outside some bounded interval).
There are straightforward mappings from both locally integrable functions and
probability distributions to corresponding distributions, as discussed below. However,
not all distributions can be formed in this manner.
Suppose that
This integral is a real number which linearly and continuously depends on . This
suggests the requirement that a distribution should be linear and continuous over the
space of test functions D(R), which completes the definition. In a conventional abuse of
notation, f may be used to represent both the original function f and the distribution Tf
derived from it.
Such distributions may be multiplied with real numbers and can be added together, so
they form a real vector space. In general it is not possible to define a multiplication for
distributions, but distributions may be multiplied with infinitely differentiable functions.
It's desirable to choose a definition for the derivative of a distribution which, at least for
distributions derived from locally integrable functions, has the property that (T f)' = Tf '.
If is a test function, we can show that
130
using integration by parts and noting that , since φ is zero
outside of a bounded set. This suggests that if S is a distribution, we should define its
derivative S' by
It turns out that this is the proper definition; it extends the ordinary definition of
derivative, every distribution becomes infinitely differentiable and the usual properties
of derivatives hold.
Example: Recall that the Dirac delta (so-called Dirac delta function) is the distribution
defined by
This latter distribution is our first example of a distribution which is derived from
neither a function nor a probability distribution.
There are several common conventions for defining the Fourier transform = F(f) of
an integrable function ƒ : R → C (Kaiser 1994). This article will use the definition
131
When the independent variable x represents time (with SI unit of seconds), the transform
variable ξ represents frequency (in hertz).
Under suitable conditions, ƒ can be reconstructed from by the inverse transform:
History
Joseph Fourier presented what is now called the Fourier integral theorem in his treatise
Théorie analytique de la chaleur in the form:[8]
Cauchy pointed out that in some circumstances the order of integration in this result
was significant.[12][13]
As justified using the theory of distributions, the Cauchy equation can be rearranged to
resemble Fourier's original formulation and expose the δ-function as:
132
where the δ-function is expressed as:
From a purely mathematical viewpoint, the Dirac delta is not strictly a function, because
any extended-real function that is equal to zero everywhere but a single point must have
total integral zero.[7] The delta function only makes sense as a mathematical object when
it appears inside an integral. While from this perspective the Dirac delta can usually be
manipulated as though it were a function, formally it must be defined as a distribution
that is also a measure. In many applications, the Dirac delta is regarded as a kind of
limit (a weak limit) of a sequence of functions having a tall spike at the origin. The
approximating functions of the sequence are thus "approximate" or "nascent" delta
functions.
where x and ξ are n-dimensional vectors, and x · ξ is the dot product of the vectors.
The dot product is sometimes written as .
Here we assume f(x), g(x), and h(x) are integrable functions, are Lebesgue-measurable
on the real line, and satisfy:
133
Basic properties
The Fourier transform has the following basic properties: (Pinsky 2002).
Linearity
Translation
Modulation
Scaling
The case a = −1 leads to the time-reversal property, which states: if h(x) = ƒ(−x),
Then .
Conjugation
If , then
Duality
If then
Convolution
If , then
The convolution of ƒ and g is written ƒ∗g, using an asterisk or star. It is defined as the
integral of the product of the two functions after one is reversed and shifted. As such, it
is a particular kind of integral transform:
134
(commutativity)
While the symbol t is used above, it need not represent the time domain. But in that
context, the convolution formula can be described as a weighted average of the function
ƒ(τ) at the moment t where the weighting is given by g(−τ) simply shifted by amount t.
As t changes, the weighting function emphasizes different parts of the input function.
More generally, if f and g are complex-valued functions on Rd, then their convolution
may be defined as the integral:
Commutativity
Associativity
Distributivity
Multiplicative identity
No algebra of functions possesses an identity for the convolution. The lack of identity is
typically not a major inconvenience, since most collections of functions on which the
convolution is performed can be convolved with a delta distribution or, at the very least
(as is the case of L1) admit approximations to the identity. The linear space of
compactly supported distributions does, however, admit an identity under the
convolution. Specifically,
135
where δ is the delta distribution.
Inverse element
Some distributions have an inverse element for the convolution, S(−1), which is defined
by
The set of invertible distributions forms an abelian group under the convolution.
Complex conjugation
Integration
If ƒ and g are integrable functions, then the integral of their convolution on the whole
space is simply obtained as the product of their integrals:
This follows from Fubini's theorem. The same result holds if ƒ and g are only assumed
to be nonnegative measurable functions, by Tonelli's theorem.
Differentiation
where d/dx is the derivative. More generally, in the case of functions of several
variables, an analogous formula holds with the partial derivative:
These identities hold under the precise condition that ƒ and g are absolutely integrable
and at least one of them has an absolutely integrable (L1) weak derivative, as a
consequence of Young's inequality. For instance, when ƒ is continuously differentiable
with compact support, and g is an arbitrary locally integrable function,
136
These identities also hold much more broadly in the sense of tempered distributions if
one of ƒ or g is a compactly supported distribution or a Schwartz function and the other
is a tempered distribution. On the other hand, two positive integrable and infinitely
differentiable functions may have a nowhere continuous convolution.
In the discrete case, the difference operator D ƒ(n) = ƒ(n + 1) − ƒ(n) satisfies an
analogous relationship:
Translation invariance
If ƒ is a Schwartz function, then τxƒ is the convolution with a translated Dirac delta
function τxƒ = ƒ∗τx δ. So translation invariance of the convolution of Schwartz functions
is a consequence of the associativity of convolution.
Convolution theorem
137
In linear time invariant (LTI) system theory, it is common to interpret g(x) as the
impulse response of an LTI system with input ƒ(x) and output h(x), since substituting
the unit impulse for ƒ(x) yields h(x) = g(x). In this case, represents the frequency
response of the system.
Conversely, if ƒ(x) can be decomposed as the product of two square integrable functions
p(x) and q(x), then the Fourier transform of ƒ(x) is given by the convolution of the
respective Fourier transforms and .
Versions of this theorem also hold for the Laplace transform, two-sided Laplace
transform, Z-transform and Mellin transform.(reference: Titchmarsh convolution
theorem).
The most common differential operator is the action of taking the derivative itself.
Common notations for taking the first derivative with respect to a variable x include:
and .
When taking higher, nth order derivatives, the operator may also be written:
or
The D notation's use and creation is credited to Oliver Heaviside, who considered
differential operators of the form
One of the most frequently seen differential operators is the Laplacian operator, defined
by
138
Another differential operator is the Θ operator, or theta operator, defined by
This is sometimes also called the homogeneity operator, because its eigenfunctions are
the monomials in z:
Del is used to calculate the gradient, curl, divergence, and laplacian of various objects.
Adjoint of an operator
See also: Hermitian adjoint
where the notation is used for the scalar product or inner product. This
definition therefore depends on the definition of the scalar product.
In the functional space of square integrable functions, the scalar product is defined by
If one moreover adds the condition that f or g vanishes for and , one
can also define the adjoint of T by
139
This formula does not explicitly depend on the definition of the scalar product. It is
therefore sometimes chosen as a definition of the adjoint operator. When is defined
according to this formula, it is called the formal adjoint of T.
Elliptic operator
From Wikipedia, the free encyclopedia
Elliptic operators are typical of potential theory, and they appear frequently in
electrostatics and continuum mechanics. Elliptic regularity implies that their solutions
tend to be smooth functions (if the coefficients in the operator are smooth). Steady-state
solutions to hyperbolic and parabolic equations generally solve elliptic equations.
Formal definitions:
In many applications, this condition is not strong enough, and instead a UNIFORM
ELLIPTICITY CONDITION may be imposed for operators of degree m = 2k:
where C is a positive constant. Note that ellipticity only depends on the highest-order
terms.
A nonlinear operator
140
is elliptic if its first-order Taylor expansion with respect to u and its derivatives about
any point is a linear elliptic operator.
-Example:
-Another example:
Given a matrix-valued function A(x) which is symmetric and positive definite for every x, having
components aij , the operator
is elliptic. This is the most general form of a second-order divergence form linear elliptic
differential operator. The Laplace operator is obtained by taking A = I. These operators also
occur in electrostatics in polarized media.
A similar nonlinear operator occurs in glacier mechanics. The stress tensor of ice, according to
Glen's flow law, is given by
for some constant B. The velocity of an ice sheet in steady state will then solve the nonlinear
elliptic system
141
where ρ is the ice density, g is the gravitational acceleration vector, p is the pressure and Q is a
forcing term.
This situation is ultimately unsatisfactory, as the weak solution u might not have enough
derivatives for the expression Lu to even make sense.
Any differential operator exhibiting this property is called a hypoelliptic operator; thus,
every elliptic operator is hypoelliptic. The property also means that every fundamental
solution of an elliptic operator is infinitely differentiable in any neighborhood not
containing 0.
General definition
for all and all . It is important to note that the definition of ellipticity in the
previous part of the article is strong ellipticity. Here is an inner product. Notice
that the are covector fields or one-forms, but the are elements of the vector bundle
upon which acts.
The quintessential example of a (strongly) elliptic operator is the Laplacian (or its
negative, depending upon convention). It is not hard to see that needs to be of even
order for strong ellipticity to even be an option. Otherwise, just consider plugging in
142
both and its negative. On the other hand, a weakly elliptic first-order operator, such as
the Dirac operator can square to become a strongly elliptic operator, such as the
Laplacian. The composition of weakly elliptic operators is weakly elliptic.
Weak ellipticity is nevertheless strong enough for the Fredholm alternative, Schauder
estimates, and the Atiyah–Singer index theorem. On the other hand, we need strong
ellipticity for the maximum principle, and to guarantee that the eigenvalues are discrete,
and their only limit point is infinity.
143
144
A PARABOLIC PARTIAL DIFFERENTIAL EQUATION is a type of second-order partial
differential equation (PDE), describing a wide family of problems in science including
heat diffusion, ocean acoustic propagation, in physical or mathematical systems with a
time variable, and which behave essentially like heat diffusing through a solid.
This equation says roughly that the temperature at a given time and point will rise or fall
at a rate proportional to the difference between the temperature at that point and the
average temperature near that point. The quantity measures how far off the
temperature is from satisfying the mean value property of harmonic functions.
where is a second order elliptic operator (implying must be positive also). Such a
system can be hidden in an equation of the form
Under broad assumptions, parabolic PDEs as given above have solutions for all x,y and
t>0. An equation of the form is considered to be parabolic if L is a
(possibly nonlinear) function of u and its first and second derivatives, with some further
conditions on L. With such a nonlinear parabolic differential equation, solutions exist
for a short time but may explode in a singularity in a finite amount of time. Hence, the
difficulty is in determining solutions for all time, or more generally studying the
singularities that arise. This is in general quite difficult, as in the solution of the
Poincaré conjecture via Ricci flow.
145
A HYPERBOLIC PARTIAL DIFFERENTIAL EQUATION of order n is a partial differential
equation (PDE) that, roughly speaking, has a well-posed initial value problem for the
first n−1 derivatives. More precisely, the Cauchy problem can be locally solved for
arbitrary initial data along any non-characteristic hypersurface. Many of the equations
of mechanics are hyperbolic, and so the study of hyperbolic equations is of substantial
contemporary interest. The model hyperbolic equation is the wave equation. In one
spatial dimension, this is
The equation has the property that, if u and its first time derivative are arbitrarily
specified initial data on the initial line t = 0 (with sufficient smoothness properties), then
there exists a solution for all time.
Definition:
Examples:
with
146
can be transformed to the wave equation, apart from lower order terms which are
inessential for the qualitative understanding of the equation.[2] This definition is
analogous to the definition of a planar hyperbola.
Consider the following system of first order partial differential equations for
unknown functions , , where
If the matrix has distinct real eigenvalues, it follows that it's diagonalizable. In this
case the system is called strictly hyperbolic.
147
A GREEN'S FUNCTION, named after the British mathematician George Green, who first
developed the concept in the 1830s is a type of function used to solve inhomogeneous
differential equations subject to specific initial conditions or boundary conditions. In the
modern study of linear partial differential equations, Green's functions are studied largely
from the point of view of fundamental solutions instead.
(1)
where is the Dirac delta function. This property of a Green's function can be exploited
to solve differential equations of the form
(2)
If the kernel of L is non-trivial, then the Green's function is not unique. However, in
practice, some combination of symmetry, boundary conditions and/or other externally
imposed criteria will give a unique Green's function. Also, Green's functions in general
are distributions, not necessarily proper functions.
Green's functions are also a useful tool in solving wave equations, diffusion equations,
and in quantum mechanics, where the Green's function of the Hamiltonian is a key
concept, with important links to the concept of density of states. As a side note, the
Green's function as used in physics is usually defined with the opposite sign; that is,
This definition does not significantly change any of the properties of the Green's
function.
If the operator is translation invariant, that is when L has constant coefficients with
respect to x, then the Green's function can be taken to be a convolution operator, that is,
In this case, the Green's function is the same as the impulse response of linear time-
invariant system theory.
Motivation
Loosely speaking, if such a function G can be found for the operator L, then if we
multiply the equation (1) for the Green's function by f(s), and then perform an
integration in the s variable, we obtain;
148
The right hand side is now given by the equation (2) to be equal to L u(x), thus:
Because the operator L = L(x) is linear and acts on the variable x alone (not on the
variable of integration s), we can take the operator L outside of the integration on the
right hand side, obtaining;
(3)
Thus, we can obtain the function u(x) through knowledge of the Green's function in
equation (1), and the source term on the right hand side in equation (2). This process
relies upon the linearity of the operator L.
In other words, the solution of equation (2), u(x), can be determined by the integration
given in equation (3). Although f(x) is known, this integration cannot be performed
unless G is also known. The problem now lies in finding the Green's function G that
satisfies equation (1). For this reason, the Green's function is also sometimes called the
fundamental solution associated to the operator L.
Not every operator L admits a Green's function. A Green's function can also be thought
of as a right inverse of L. Aside from the difficulties of finding a Green's function for a
particular operator, the integral in equation (3), may be quite difficult to evaluate.
However the method gives a theoretically exact result.
Framework
Let f(x) be a continuous function in [0,l]. We shall also suppose that the problem
is regular (i.e., only the trivial solution exists for the homogeneous problem).
Theorem
and it is given by
Eigenvalue expansions
150
If a differential operator L admits a set of eigenvectors (i.e., a set of functions
and scalars such that ) that is complete, then it is possible to
construct a Green's function from these eigenvectors and eigenvalues.
Complete means that the set of functions satisfies the following completeness
relation:
Applying the operator L to each side of this equation results in the completeness
relation, which was assumed true.
The general study of the Green's function written in the above form, and its relationship
to the function spaces formed by the eigenvectors, is known as Fredholm theory.
Green's functions for linear differential operators involving the Laplacian may be
readily put to use using the second of Green's identities.
To derive Green's theorem, begin with the divergence theorem (otherwise known as
Gauss's theorem):
151
Suppose that the linear differential operator L is the Laplacian, , and that there is a
Green's function G for the Laplacian. The defining property of the Green's function still
holds:
Suppose the problem is to solve for inside the region. Then the integral
reduces to simply due to the defining property of the Dirac delta function and
we have:
This form expresses the well-known property of harmonic functions that if the value or
normal derivative is known on a bounding surface, then the value of the function inside
the volume is known everywhere.
152
If the problem is to solve a Dirichlet boundary value problem, the Green's function
should be chosen such that vanishes when either x or x' is on the bounding
surface.Thus only one of the two terms in the surface integral remains. If the problem is
to solve a Neumann boundary value problem, the Green's function is chosen such that
its normal derivative vanishes on the bounding surface, as it would seems to be the most
logical choice. (See Jackson J.D. classical electrodynamics, page 39). However,
application of Gauss's theorem to the differential equation defining the Green's function
yields
where is the average value of the potential on the surface. This number is not
known in general, but is often unimportant, as the goal is often to obtain the electric
field given by the gradient of the potential, rather than the potential itself.
With no boundary conditions, the Green's function for the Laplacian (Green's function
for the three-variable Laplace equation) is:
Supposing that the bounding surface goes out to infinity, and plugging in this
expression for the Green's function, this gives the familiar expression for electric
potential in terms of electric charge density (in the CGS unit system) as
Example
153
find the Green's function.
First step: The Green's function for the linear operator at hand is defined as the solution
to
If , then the delta function gives zero, and the general solution is
One can also ensure proper discontinuity in the first derivative by integrating the
154
So the Green's function for this problem is:
EIGENVALUES and EIGENVECTORS have many applications in both pure and applied
mathematics. They are used in matrix factorization, in quantum mechanics, and in many other
areas.
History
Eigenvalues are often introduced in the context of linear algebra or matrix theory.
Historically, however, they arose in the study of quadratic forms and differential
equations.
Euler studied the rotational motion of a rigid body and discovered the importance of the
principal axes. Lagrange realized that the principal axes are the eigenvectors of the
inertia matrix.[11] In the early 19th century, Cauchy saw how their work could be used to
classify the quadric surfaces, and generalized it to arbitrary dimensions. [12] Cauchy also
coined the term racine caractéristique (characteristic root) for what is now called
eigenvalue; his term survives in characteristic equation.[13]
Fourier used the work of Laplace and Lagrange to solve the heat equation by separation
of variables in his famous 1822 book Théorie analytique de la chaleur. Sturm
developed Fourier's ideas further and brought them to the attention of Cauchy, who
combined them with his own ideas and arrived at the fact that real symmetric matrices
have real eigenvalues. This was extended by Hermite in 1855 to what are now called
Hermitian matrices. Around the same time, Brioschi proved that the eigenvalues of
orthogonal matrices lie on the unit circle, and Clebsch found the corresponding result
for skew-symmetric matrices. Finally, Weierstrass clarified an important aspect in the
stability theory started by Laplace by realizing that defective matrices can cause
instability.
In the meantime, Liouville studied eigenvalue problems similar to those of Sturm; the
discipline that grew out of their work is now called Sturm–Liouville theory. Schwarz
studied the first eigenvalue of Laplace's equation on general domains towards the end of
the 19th century, while Poincaré studied Poisson's equation a few years later.
At the start of the 20th century, Hilbert studied the eigenvalues of integral operators by
viewing the operators as infinite matrices. He was the first to use the German word
eigen to denote eigenvalues and eigenvectors in 1904, though he may have been
following a related usage by Helmholtz. For some time, the standard term in English
was "proper value", but the more distinctive term "eigenvalue" is standard today. [18]
155
The first numerical algorithm for computing eigenvalues and eigenvectors appeared in
1929, when Von Mises published the power method. One of the most popular methods
today, the QR algorithm, was proposed independently by John G.F. Francis and Vera
Kublanovskaya in 1961.
The EIGENVECTORS of a square matrix are the non-zero vectors that, after being
multiplied by the matrix, remain parallel to the original vector. For each eigenvector, the
corresponding eigenvalue is the factor by which the eigenvector is scaled when
multiplied by the matrix. The prefix eigen- is adopted from the German word "eigen"
for "own" in the sense of a characteristic description. The eigenvectors are sometimes
also called characteristic vectors. Similarly, the eigenvalues are also known as
characteristic values.
These ideas are often extended to more general situations, where scalars are elements of
any field, vectors are elements of any vector space, and linear transformations may or
may not be represented by matrix multiplication. For example, instead of real numbers,
scalars may be complex numbers; instead of arrows, vectors may be functions or
frequencies; instead of matrix multiplication, linear transformations may be operators
such as the derivative from calculus. These are only a few of countless examples where
eigenvectors and eigenvalues are important.
In such cases, the concept of direction loses its ordinary meaning, and is given an
abstract definition. Even so, if that abstract direction is unchanged by a given linear
transformation, the prefix "eigen" is used, as in eigenfunction, eigenmode, eigenface,
eigenstate, and eigenfrequency.
Eigenvalues and eigenvectors have many applications in both pure and applied
mathematics. They are used in matrix factorization, in quantum mechanics, and in many
other areas.
The eigenvalue equation for linear differential operators is then a set of one or more
differential equations. The eigenvectors are commonly called eigenfunctions. The
simplest case is the eigenvalue equation for differentiation of a real valued function by a
single real variable. We seek a function (equivalent to an infinite-dimensional vector)
that, when differentiated, yields a constant times the original function. In this case, the
eigenvalue equation becomes the linear differential equation
156
Here λ is the eigenvalue associated with the function, f(x). This eigenvalue equation has
a solution for any value of λ. If λ is zero, the solution is
If we expand our horizons to complex valued functions, the value of λ can be any
complex number. The spectrum of d/dt is therefore the whole complex plane. This is an
example of a continuous spectrum.
Waves on a string: The shape of a standing wave in a string fixed at its boundaries is an
example of an eigenfunction of a differential operator. The admittable eigenvalues are
governed by the length of the string and determine the frequency of oscillation.
The displacement, , of a stressed rope fixed at both ends, like the vibrating
strings of a string instrument, satisfies the wave equation
which is a linear partial differential equation, where c is the constant wave speed. The
normal method of solving such an equation is separation of variables. If we assume that
h can be written as the product of the form X(x)T(t), we can form a pair of ordinary
differential equations:
and
Each of these is an eigenvalue equation (the unfamiliar form of the eigenvalue is chosen
merely for convenience). For any values of the eigenvalues, the eigenfunctions are
given by
and
If we impose boundary conditions (that the ends of the string are fixed with X(x) = 0 at x
= 0 and x = L, for example) we can constrain the eigenvalues. For those boundary
conditions, we find
and
157
Thus, the constant is constrained to take one of the values , where n is
any integer. Thus the clamped string supports a family of standing waves of the form
From the point of view of our musical instrument, the frequency is the frequency
of the nth harmonic, which is called the (n-1)st overtone.
Applications
Schrödinger equation
The wavefunctions associated with the bound states of an electron in a hydrogen atom can be
seen as the eigenvectors of the hydrogen atom Hamiltonian as well as of the angular
momentum operator. They are associated with eigenvalues interpreted as their energies
(increasing downward: n=1,2,3,...) and angular momentum (increasing across: s, p, d,...). The
illustration shows the square of the absolute value of the wavefunctions. Brighter areas
correspond to higher probability density for a position measurement. The center of each figure
is the atomic nucleus, a proton.
However, in the case where one is interested only in the bound state solutions of the
Schrödinger equation, one looks for within the space of square integrable functions.
Since this space is a Hilbert space with a well-defined scalar product, one can introduce
a basis set in which and H can be represented as a one-dimensional array and a
matrix respectively. This allows one to represent the Schrödinger equation in a matrix
form.
Bra-ket notation is often used in this context. A vector, which represents a state of the
system, in the Hilbert space of square integrable functions is represented by . In
this notation, the Schrödinger equation is:
158
where is an eigenstate of H. It is a self adjoint operator, the infinite dimensional
analog of Hermitian matrices (see Observable). As in the matrix case, in the equation
above is understood to be the vector obtained by application of the
transformation H to .
or
where is the eigenvalue and is the angular frequency. Note that the principal
vibration modes are different from the principal compliance modes, which are the
eigenvectors of k alone. Furthermore, damped vibration, governed by
This can be reduced to a generalized eigenvalue problem by clever algebra at the cost of
solving a larger system.
(2)
159
PODE SER EXPRESSA POR
(3)
is an iterated partial derivative, where ∂ j means differentiation with respect to the j-th
, where
160
one obtains the operator written as a composition of a Fourier transform, a simple
multiplication by the polynomial function and an inverse Fourier transform:
(1)
we (formally) apply the Fourier transform on both sides and obtain the algebraic
equation
If the symbol P(ξ) is never zero when ξ ∈ Rn, then it is possible to divide by P(ξ):
The last assumption can be weakened by using the theory of distributions. The first two
assumptions can be weakened as follows.
This is similar to formula (1), except that 1/P(ξ) is not a polynomial function, but a
function of a more general kind.
161
Formal Definition of Pseudo-Differential Operators
(2)
where the symbol P(x,ξ) in the integrand belongs to a certain symbol class. For instance,
if P(x,ξ) is an infinitely differentiable function on Rn × Rn with the property
for all x,ξ ∈Rn, all multiindices α,β. some constants Cα, β and some real number m, then
P belongs to the symbol class of Hörmander. The corresponding operator P(x,D)
is called a pseudo-differential operator of order m and belongs to the class
Properties
Linear differential operators of order m with smooth bounded coefficients are pseudo-
differential operators of order m. The composition PQ of two pseudo-differential
operators P, Q is again a pseudo-differential operator and the symbol of PQ can be
calculated by using the symbols of P and Q. The adjoint and transpose of a pseudo-
differential operator is a pseudo-differential operator.
Differential operators are local in the sense that one only needs the value of a function
in a neighbourhood of a point to determine the effect of the operator. Pseudo-differential
operators are pseudo-local, which means informally that when applied to a distribution
they do not create a singularity at points where the distribution was already smooth.
162
FOURIER INTEGRAL OPERATORS have become an important tool in the theory of partial
differential equations. This class of operators contains differential operators as well as
classical integral operators as special cases.
One motivation for the study of Fourier integral operators is the solution operator for
the initial value problem for the wave operator. Indeed, consider the following problem:
and
L. Hörmander Fourier integral operators, Acta Math. 127 (1971), 79–183. doi
10.1007/BF02392052, http://www.springerlink.com/content/t202410l4v37r13m/fulltext.pdf
163
OBSERVAÇÃO: NA DÉCADA DE 1970, O PROFESSOR ANTONIO GILIOLI CONSULTOU SEU
ORIENTADOR, FRANÇOIS TREVES, SOBRE COMO INTERPRETAVA AS GENERALIZAÇÕES
MATEMÁTICAS QUE VINHAM SENDO PUBLICADAS, E OBTEVE COMO RESPOSTA QUE ELE
ABORDAVA EQUAÇÕES DIFERENCIAIS E ERA RECEPTIVO A TODA A MATEMÁTICA
PERTINENTE. POSTERIORMENTE, FOMOS INFORMADOS POR SEU OUTRO ORIENTADO,
PROFESSOR PAULO DOMINGOS CORDARO DE QUE, APÓS A ABORDAGEM DOS OPERADORES
DE TIPO PRINCIPAL, ELE ESTARIA PROPENSO A SE DEDICAR AOS SISTEMAS DE EQUAÇÕES
DIFERENCIAIS.
SEGUNDO ALGUNS ESPECIALISTAS, NA DÉCADA DE 1970, A TEORIA DAS EDP LINEARES JÁ
TERIA ALCANÇADO UM NÍVEL DE ELABORAÇÃO ELEVADO, NO QUAL OS GANHOS EM
GENERALIDADE, LIMITADOS PELOS CONTRA-EXEMPLOS EXISTENTES, NEM SEMPRE ERAM
INSPIRADOS EM APLICAÇÕES RELEVANTES, E VINHAM, AINDA, ACOMPANHADOS DE MAIOR
COMPLEXIDADE. AO MESMO TEMPO SURGIRAM PUBLICAÇÕES QUE PODEM SER
CLASSIFICADAS COMO TEOREMAS DE REGULARIDADE DAS SOLUÇÕES (PROPRIEDADES DAS
SOLUÇÕES DE CLASSES ESPECIAIS DE EQUAÇÕES). OS PROFESSORES DJAIRO GUEDES DE
FIGUEIREDO E LOUIS NIRENBERG FORAM CONSULTADOS E EXPRESSARAM O
ENTENDIMENTO DE QUE A ABORDAGEM DAS EDP NÃO LINEARES SE CONSTITUIA NUM
RAMO MAIS PROMISSOR OU COM MELHORES PERSPECTIVAS PARA O NOSSO AMBIENTE, DO
QUE O DAS EDP LINEARES, JÁ ELABORADO.
164
NONLINEAR PARTIAL DIFFERENTIAL EQUATIONS describe many different physical
systems, ranging from gravitation to fluid dynamics, and have been used in mathematics
to solve problems such as the Poincaré conjecture and the Calabi conjecture. They are
difficult to study: there are almost no general techniques that work for all such
equations, and usually each individual equation has to be studied as a separate problem
A fundamental question for any PDE is the existence and uniqueness of a solution for
given boundary conditions. For nonlinear equations these questions are in general very
hard: for example, the hardest part of Yau's solution of the Calabi conjecture was the
proof of existence for a Monge–Ampere equation.
The basic questions about singularities (their formation, propagation, and removal, and
regularity of solutions) are the same as for linear PDE, but as usual much harder to
study. In the linear case one can just use spaces of distributions, but nonlinear PDEs are
not usually defined on arbitrary distributions, so one replaces spaces of distributions by
refinements such as Sobolev spaces.
An example of singularity formation is given by the Ricci flow: Hamilton showed that
while short time solutions exist, singularities will usually form after a finite time.
Perelman's solution of the Poincaré conjecture depended on a deep study of these
singularities, where he showed how to continue the solution past the singularities.
Linear approximation
A–F
Benjamin– Fluid
1+1
Bona–Mahony mechanics
internal
Benjamin-Ono 1+1 waves in
deep water
Born-Infeld 1+1
165
Fluid
Boussinesq 1+1
mechanics
Thin viscous
Buckmaster 1+1 fluid sheet
flow
Fluid
Burgers 1+1
mechanics
Cahn–Hilliard Phase
Any
equation separation
Calabi–Yau
Calabi flow Any
manifolds
Camassa–
1+1 Peakons
Holm
Carleman 1+1
Cauchy Momentum
any
momentum transport
Caudrey–
Dodd–
Gibbon– 1+1 Same as (rescaled) Sawada–Kotera
Sawada–
Kotera
Clairaut Differential
any
equation geometry
Complex
Calabi
Monge– Any lower order terms conjecture
Ampère
Degasperis–
1+1 Peakons
Procesi
166
Dispersive
1+1 ,
long wave
Drinfel'd–
Sokolov– 1+1
Wilson
Eckhaus Integrable
1+1
equation systems
Eikonal
any optics
equation
Ernst
2
equation
Euler non-viscous
1+3
equations fluids
Fisher's Gene
1+1
equation propagation
Fitzhugh-
1+1
Nagumo
G–K
Di
Name Equation Applications
m
Gardner
1+1
equation
167
Garnier isomonodromic
equation deformations
Gauss–
surfaces
Codazzi
Ginzburg– Superconductivit
1+3
Landau y
Gross–
1+1
Neveu
Gross– Bose–Einstein
1+n
Pitaevskii condensate
Hartree
Any
equation
where .
Hasegawa– Turbulence in
1+3
Mima plasma
Heisenberg
ferromagne 1+1 Magnetism
t
Hirota
1+1
equation
Hirota– ,
1+1
Satsuma
Hunter–
1+1 Liquid crystals
Saxton
Ishimori Integrable
1+2
equation systems
168
–Petviashvili waves
,
von Karman 2
Kaup 1+1
Kaup–
Integrable
Kupershmid 1+1
systems
t
Klein–
Gordon– any ,
Maxwell
Klein–
Gordon any
(nonlinear)
Klein–
Gordon–
Zakharov
Khokhlov–
Zabolotskay 1+2
a
KdV
(generalized 1+1
)
KdV
1+1
(modified)
,
KdV (super) 1+1
There are more minor variations listed in the article on KdV equations.
Kuramoto–
1+n
Sivashinsky
169
L–R
Di Applicati
Name Equation
m ons
Landau– Magnetic
1+
Lifshitz field in
n
model solids
Lin-Tsien 1+
equation 2
an
Liouville
y
Minimal minimal
3
surface surfaces
Molenbro
2
eck
Monge– an
Ampère y lower order terms
Navier–
Stokes
1+
(and its Fluid flow
3
derivation + mass conservation:
) + an equation of state to relate p and ρ, e.g. for an
incompressible flow:
Nonlinear optics,
1+
Schröding water
1
er (cubic) waves
Nonlinear
Schröding optics,
1+
er water
1
(derivativ waves
e)
170
Novikov–
1+
Veselov see Veselov–Novikov equation below
2
equation
atmosph
Omega 1+
eric
equation 3
physics
Plateau 2
Pohlmeye
r–Lund– 2
Regge
Porous 1+
diffusion
medium n
1+ boundary
Prandtl
2 , layer
Atmosph
Primitive 1+
eric
equations 3
models
S–Z, α–ω
Di
Name Equation Applications
m
Rayleigh 2
An Poincaré
Ricci flow
y conjecture
Variably-
saturated
Richards 1+
flow in
equation 3
porous
media
Sawada– 1+
171
Kotera 1
isomonodro
Schlesing An mic
er y deformation
s
Seiberg–
Seiberg– 1+ Witten
Witten 3 invariants,
QFT
Shallow 1+ shallow
water 2 water waves
Sine– 1+
Solitons, QFT
Gordon 1
Sinh– 1+
Solitons, QFT
Gordon 1
Sinh– 1+
Poisson n
Swift–
an pattern
Hohenbe
y forming
rg
Three-
1+ Integrable
wave
n systems
equation
Thomas
2
equation
172
Toda an
lattice y
Veselov–
1+ , shallow
Novikov
2 , water waves
equation
Wadati–
Konno–
1+
Ichikawa
1
–
Schimizu
WDVV Topological
An
equation field theory,
y
s QFT
WZW 1+
QFT
model 1
Witham phase
equation averaging
Differential
Yamabe n
geometry
Yang–
Mills
An Gauge
equation
y theory, QFT
(source-
free)
Yang–
Mills
Instantons,
(self-
4 Donaldson
dual/anti
theory, QFT
-self-
dual)
Meson-
Yukawa 1+ nucleon
equation n interactions,
QFT
173
Zakharov 1+ Langmuir
system 3 waves
Zakharov
– 1+ Acoustic
Schulma 3 waves
n
Zoomero 1+
Solitons
n 1
φ4 1+
QFT
equation 1
Harmonic
1+ maps,
σ-model
1 integrable
systems, QFT
These lectures are meant as an informal introduction to some of the techniques used in
proving existence of solutions of nonlinear problems of the form
F(u) = y. (1)
Here F is a continuous (and usually smooth) mapping from one topological space X to
another Y, and the spaces are usually infinite dimensional. The model to keep in mind is
one in which these spaces are function spaces defined in domains on finite-dimensional
manifolds, and F is a system of nonlinear partial differential operators-or integral
operators.
A number of special topics will be presented-in three parts:
174
I. Global methods: homotopy, in particular topological degree theory, and extensions.
Applications to nonlinear boundary value problems.
II. Variational methods, in which a solution is a stationary point of some functional.
Applications.
III.Local study, perturbation about a solution.
An important analytic aspect of all these problems is that of finding a priori estimates
for the solutions. How one does that varies from problem to problem and I will barely
touch on these technical aspects. I will try, rather, to avoid technicalities and stress the
topological and variational ideas.
The lectures are not addressed to the experts in these fields-for them there will be little
new. They are given with the hope of attracting others to the subject. Up to now, the
topological and abstract ideas used are rather primitive, and I am confident that there
will be enormous further development-
involving more and more sophisticated topology. Here is a more specific list of the
topics treated.
II. 1 The Palais-Smale condition (PS) and an elliptic boundary Value problem.
11.2 The mountain pass lemma. Solution of the boundary value problem.
11.3 Generalizations and variants of the mountain pass lemma.
11.4 A theorem of P. Rabinowitz on periodic solutions for a Hamiltonian system.
11.5 Periodic solutions of a nonlinear string equation.
…..
175
VI-A FUNDAMENTAÇÃO DA ANÁLISE, AS GENERALIZAÇÕES E
AS DECORRENTES QUESTÕES DE FUNDAMENTAÇÃO
65- EM MEADOS DO SÉCULO XVIII, HOUVE DIVERGÊNCIA ENTRE d’ALEMBERT, EULER E
BERNOULLI, SOBRE A SOLUÇÃO DA EQUAÇÃO DAS ONDAS , EM VIRTUDE DO USO DE
DISTINTAS CONCEPÇÕES PARA O CONCEITO DE FUNÇÃO, À ÉPOCA, AINDA NÃO ASSENTADO
(NOTAS HISTÓRICAS-DJAIRO GUEDES DE FIGUEIREDO/ CARSLAW);
-EM 1811, FOURIER ENUNCIOU QUE TODA FUNÇÃO PODERIA SER EXPRESSA POR UMA SÉRIE
TRIGONOMÉTRICA, COM OS CHAMADOS COEFICIENTES DE FOURIER, OU SEJA, POR SUA SÉRIE
DE FOURIER, O QUE INDUZIU UMA REVISÃO DOS CONCEITOS DA ANÁLISE, INICIADA POR
CAUCHY (CUJAS NOTAS DO CURSO DE CÁLCULO DE 1823, SE ENCONTRAM NA BIBLIOTECA DO
IME), BOLZANO (QUE APRESENTOU EM 1834, O 1º EXEMPLO DE UMA FUNÇÃO CONTÍNUA,
QUE NÃO É DERIVÁVEL EM NENHUM PONTO), RIEMANN E OUTROS;
PARA A INVESTIGAÇÃO DA CONVERGÊNCIA DA SÉRIE DE FOURIER DE UMA FUNÇÃO f PARA f ,
E NA TENTATIVA DE AMPLIAR A CLASSE DE FUNÇÕES INTEGRÁVEIS, RIEMANN LANÇOU MÃO
DE UMA “DEFINIÇÃO OPERACIONAL”, CHAMANDO DE INTEGRÁVEIS, À CLASSE DE TODAS AS
FUNÇÕES LIMITADAS PARA AS QUAIS O PROCESSO DE INTEGRAÇÃO PUDESSE SER
EFETIVADO, OU SEJA, PARA AS QUAIS, A INTEGRAL SUPERIOR FOSSE IGUAL À INTEGRAL
INFERIOR, PARA, EM SEGUIDA, CARACTERIZÁ-LAS ATRAVÉS DE ALGUMA PROPRIEDADE MAIS
OPERACIONALMENTE VERIFICÁVEL, OBTENDO QUE: “UMA FUNÇÃO LIMITADA, DEFINIDA
NUM INTERVALO, É INTEGRÁVEL, SE, E SOMENTE SE, O CONJUNTO DOS PONTOS DE
DESCONTINUIDADE MEDIDA (DE LEBESGUE) NULA” ;
176
EMBORA SEM INTERIOR (NÃO CONTENDO NENHUM INTERVALO), TEM MEDIDA POSITIVA;
(COUNTEREXAMPLES IN ANALYSIS- B.GELBAUM/J.OLMSTED);
177
AS LIMITAÇÕES DOS SISTEMAS FORMAIS
68- NA INTENÇÃO DE ESCLARECER AS DIVERSAS QUESTÕES ENVOLVIDAS, TENTOU-SE
APRESENTAR, PELO MÉTODO AXIOMÁTICO, (PRIMEIRO SEPARADAMENTE PARA DEPOIS
SEREM COMPOSTAS), TEORIAS (CONTENDO PRINCÍPIOS E SUAS CONSEQUÊNCIAS), PARA CADA
UM DOS COMPONENTES FUNDAMENTAIS DAS TEORIAS MATEMÁTICAS: METALINGUAGEM,
LÓGICA SUBJACENTE, LINGUAGEM FORMAL, TEORIAS DE FUNDAMENTAÇÃO, TEORIAS
MATEMÁTICAS (PROPRIAMENTE DITAS), E ESTRUTURAS ABSTRATAS, ALÉM DE TEORIA DE
MODELOS E INTERPRETAÇÕES DAS TEORIAS FORMAIS, TENDO A FORMALIDADE COMO
CRITÉRIO DE VALIDAÇÃO:
178
ENVOLVEM TODA A TEORIA (COMO A SUA CONSISTÊNCIA), QUE SÃO ENUNCIADOS DA
METATEORIA, PODEM SER TRANSCRITAS NA TEORIA, MAS A DECISÃO DE QUE SÃO
VERDADEIRAS OU FALSAS (COM RECURSOS SOMENTE DA TEORIA), PODE GERAR
CONTRADIÇÃO;
179
NOS ESPAÇOS VETORIAIS ( V ), ALÉM DA ESTRUTURA DE GRUPO COMUTATIVO (NO QUAL,
ALÉM DAS PROPRIEDADES CITADAS, A OPERAÇÃO TAMBÉM É COMUTATIVA), HÁ A OPERAÇÃO
DE PRODUTO POR ESCALAR (EM ANÁLISE, OU ) , QUE A CADA ESCALAR λ , E
VETOR x V, ASSOCIA O VETOR λx V , DE MODO QUE: FIXANDO x V,E
VARIANDO λ EM , OBTEMOS UMA “CÓPIA DE EM V” , NA DIREÇÃO DO VETOR x ,
ALÉM DE OUTRAS PROPRIEDADES OPERACIONAIS;
180
-EM 1854, G.RIEMANN APRESENTOU UM TRATAMENTO UNIFICADO PARA AS GEOMETRIAS
EUCLIDIANA E NÃO EUCLIDIANAS, NA CHAMADA GEOMETRIA RIEMANNIANA, NA QUAL, EM
TODO PONTO, O ESPAÇO TANGENTE POSSUI UMA MÉTRICA POSITIVA-DEFINIDA (MÉTRICA
RIEMANNIANA), QUE POSTERIORMENTE FOI GENERALIZADA À GEOMETRIA PSEUDO-
RIEMANNIANA, NA QUAL A MÉTRICA NÃO PRECISA SER POSITIVA-DEFINIDA, E ENGLOBA A
GEOMETRIA DO ESPAÇO-TEMPO DA RELATIVIDADE GERAL (ONDE O TRANSPORTE PARALELO
DE VETORES TANGENTES PERMITE A DERIVAÇÃO);
- EM 1913, H. WEYL APRESENTOU A DEFINIÇÃO MODERNA, PARA UMA VARIEDADE
DIFERENCIÁVEL DE DIMENSÃO 2 ; E,
- EM 1936, H.WHITNEY DEFINIU AS VARIEDADES DIFERENCIÁVEIS EM TERMOS DE CARTAS
LOCAIS, OU SEJA: UM ESPAÇO TOPOLÓGICO DE HAUSDORFF, NO QUAL TODO PONTO POSSUI
UMA VIZINHANÇA HOMEOMORFA A UM ABERTO DO Rn , PARA ALGUM n FIXO; PARA
TRANSPORTAR A ESTRUTURA DIFERENCIÁVEL DO Rn PARA A VARIEDADE, PEDE-SE QUE:
QUANDO DUAS CARTAS SE INTERCEPTAM, A COMPOSIÇÃO DE UMA DELAS COM A INVERSA
DA OUTRA (QUE É UMA APLICAÇÃO DE UM ABERTO DO Rn NO Rn ) SEJA DIFERENCIÁVEL;
181
FORAM CHAMADOS “ESPAÇOS DE BANACH” , PARA POSTERIORMENTE SEREM INTRODUZIDOS
OS ESPAÇOS VETORIAIS TOPOLÓGICOS;
para todos
para todos e ).
POSITIVA DEFINIDA ( e igual a zero se e somente se v = 0 );
NUM ESPAÇO VETORIAL COM PRODUTO INTERNO, PODE-SE DEFINIR A NORMA DE UM
VETOR v POR
182
E, EM PREFERÊNCIA À NOÇÃO DE ÂNGULO, OS CONCEITOS DE ORTOGONALIDADE
ORTONORMAL, SE , em que
δij = 1, se i = j
δij = 0, se i ≠ j
OU SEJA, SE OS VETORES TEM TODOS NORMA IGUAL A 1 , E SÃO DOIS A DOIS ORTOGONAIS;
E, QUE É UMA BASE DE V , SE TODO VETOR DE V , SE EXPRESSA DE UM ÚNICO MODO
COMO COMBINAÇÃO LINEAR FINITA DE ELEMENTOS DE B
COM A VANTAGEM DE: UMA VEZ FIXADA UMA BASE, PODERMOS, AO INVÉS DE TRABALHAR
COM OS VETORES, TRABALHAR COM OS SEUS COEFICIENTES, QUE SÃO ESCALARES;
FÓRMULA DE EULER ,
183
SUAS EQUIVALENTES FORMAS TRIGONOMÉTRICAS, FORMAM UM SISTEMA ORTOGONAL
COMPLETO DO ESPAÇO L2[0; 2 π] , DAS FUNÇÕES PERIÓDICAS, COM PERÍODO 2 π ,
DEFINIDAS NO INTERVALO [0 , 2 π ] , A VALORES COMPLEXOS (OU REAIS), REDUZINDO O
SEU ESTUDO À SEQUÊNCIA DOS SEUS COEFICIENTES NA EXPANSÃO (SÉRIE DE FOURIER), NA
NORMA DO ESPAÇO L2[0; 2 π], QUE SÃO OS COEFICIENTES DE FOURIER, E PERTENCEM AO
Sendo:
Forma Complexa:
onde:
Forma Trigonométrica:
onde:
184
73- OBSERVAÇÕES: 1) OS ESPAÇOS TOPOLÓGICOS (NOS QUAIS TODO PONTO ADMITE UM
SISTEMA DE VIZINHANÇAS, QUE FORNECEM AS PROXIMIDADES DO PONTO, POR EXEMPLO, AS
BOLAS ABERTAS DE CENTRO x E RAIO DELTA
DOS ESPAÇOS MÉTRICOS) SÃO ENTENDIDOS COMO O CONTEXTO MAIS GERAL, ONDE SE PODE
DEFINIR A NOÇÃO DE CONTINUIDADE, CONVERGÊNCIA, COMPACIDADE, ETC.;
A TÍTULO DE EXEMPLO, OS MATEMÁTICOS VINHAM TRABALHANDO COM ESPAÇOS
TOPOLÓGICOS DE HAUSDORFFF (ONDE CADA DOIS PONTOS PODEM SER “SEPARADOS” POR
VIZINHANÇAS QUE NÃO SE INTERCEPTAM), QUANDO, EM MEADOS DA DÉCADA DE 1960,
SURGIU UM EXEMPLO SIGNIFICATIVO, PROVENIENTE DA TOPOLOGIA ALGÉBRICA, DE UM
ESPAÇO TOPOLÓGICO QUE NÃO É DE HAUSDORFF, APÓS O QUE, SURGIRAM VÁRIOS
“AXIOMAS DE SEPARAÇÃO”, TOTALIZANDO MAIS DE VINTE;
2) NA CONSIDERAÇÃO DE ESPAÇOS ABSTRATOS COM DUAS ESTRUTURAS, PEDE-SE A
COMPATIBILIDADE DESSAS ESTRUTURAS, POR EXEMPLO, NOS ESPAÇOS VETORIAIS
TOPOLÓGICOS, PEDE-SE QUE AS OPERAÇÕES DE SOMA E PRODUTO POR ESCALAR SEJAM
CONTÍNUAS; TAMBÉM, COMO PROPRIEDADES BÁSICAS, É REQUERIDO QUE AS QUATRO
OPERAÇÕES E A COMPOSIÇÃO COM FUNÇÕES CONTÍNUAS, SEJAM CONTÍNUAS;
3) EM GERAL, EQUAÇÕES RELEVANTES APRESENTAM RESULTADOS ESPECÍFICOS, QUE NÃO
DERIVAM DOS RESULTADOS GERAIS VÁLIDOS NAS ESTRUTURAS ABSTRATAS (64 g);
185
RELEVANTES; E QUE
-MUITAS GENERALIZAÇÕES, AO AMPLIAREM O CONTEXTO, MODIFICAM DE MODO
ESSENCIAL OS PROBLEMAS ORIGINAIS. Por exemplo, embora a derivação para funções com
variáveis e valores complexos seja formalmente a mesma, apresenta propriedades
essencialmente distintas das satisfeitas para funções com variável real.
75-ENSINO/APRENDIZADO
186
soluções, mas o Sr, invertendo o processo, está trazendo uma solução, em busca de
um problema!” Dizia também que diante de certas questões, era necessário
“alisar o gato” (com o significado de refletir mais detidamente).
Browder: Mathematics over the last fifty years has become highly specialized,
and unfortunately the specialties are so divorced from each other that
sometimes they don’t have the appropriate degree of mutual understanding and
interaction. This is a very serious question that concerns graduate education.
Everybody is conscious of it, but nobody knows what to do about it. We are
training people following the principle enunciated by Carl Becker, who said that a
specialist is one who knows more and more about less and less until he knows
everything about nothing. Unfortunately, this is a principle which, in order to
survive, people have to practice.
187
example, the AMS meeting in the summer of the year 2000. I was rather surprised to
see this coming out so strongly in the surveys and in people’s comments, which were
not, as far as I know, related to the year 2000 meeting. We have to bring the attention
of both the mathematical community and the external publics—other scientific
disciplines, as well as the intellectual public and the general public—to the fact that
mathematics is an enterprise with enormous perspectives, in terms of its own
objectives and its impact on other forms of knowledge. There has been an enormous
growth in the practical importance of mathematics in the world, as represented by the
computer revolution and the necessity of analysis of complex systems, which
are everywhere around us. We see this in the genome project, in the mathematics
of finance, in the enormous growth in mathematical sophistication and interaction in
all the realms of theoretical physics, or even practical physics for that matter.
What we have to emphasize is that these are not just applications of known
mathematics. They are enormous growing points for mathematics. Mathematics
does no good by trying to pretend it is insulated from these things, because
these are where the vital areas of activity are. And this is not surprising. It’s a
reversion to what has been true in the past. It was just a short period of about
twenty or thirty years in which mathematicians had the delusion that
mathematics could be totally separated from other fields of human knowledge
and activity.
188
SATISFAZEM?”, EM VEZ DE SE GASTAR TANTA ENERGIA TENTANDO LECIONAR AS TEORIAS
DOS CONJUNTOS E FUNÇÕES ABSTRATAS ?
(3) Such theory may be an important means of aiming at exact theory in the
nonphysical fields of science.
Subjects like complexity, self-organization, connectionism and adaptive systems had already
been studied in the 1940s and 1950s. In fields like cybernetics, researchers like Norbert
Wiener, William Ross Ashby, John von Neumann and Heinz von Foerster examined complex
systems using mathematics. John von Neumann discovered cellular automata and self-
reproducing systems, again with only pencil and paper. Aleksandr Lyapunov and Jules Henri
Poincaré worked on the foundations of mathematical dynamical systems (stability, periodic
189
orbits of the solutions of ordinary differential equations, and chaos theory without knowing
its expression formula ) without any computer at all. At the same time Howard T. Odum, the
radiation ecologist, recognised that the study of general systems required a language that
could depict energetics, thermodynamic and kinetics at any system scale. Odum developed a
general systems, or Universal language, based on the circuit language of electronics to fulfill
this role, known as the Energy Systems Language. Between 1929-1951, Robert Maynard
Hutchins at the University of Chicago had undertaken efforts to encourage innovation and
interdisciplinary research in the social sciences, aided by the Ford Foundation with the
interdisciplinary Division of the Social Sciences established in 1931. [13] Numerous scholars had
been actively engaged in ideas before (Tectology of Alexander Bogdanov published in 1912-
1917 is a remarkable example), but in 1937 von Bertalanffy presented the general theory of
systems for a conference at the University of Chicago.
Pioneers
SYSTEM DYNAMICS was founded in the late 1950s by Jay W. Forrester of the MIT Sloan School
of Management with the establishment of the MIT System Dynamics Group. Jay W. Forrester
joined the faculty of the Sloan School at MIT in 1956, where he then developed what is now
System Dynamics. The first published article by Jay W. Forrester in the Harvard Business
Review on "Industrial Dynamics", was published in 1958. The members of the System
Dynamics Society have chosen 1957 to mark the occasion as it is the year in which the work
leading to that article, which described the dynamics of a manufacturing supply chain, was
done.
Developments
190
1990 Complex adaptive systems (CAS), John H. Holland, Murray Gell-Mann, W. Brian
Arthur
General Systems Theory - StatPac Survey Research Library 1993, David S. Walonick, Ph.D.
-Segundo Alfred Kuhn, 1974. The Logic of Social Systems. San Francisco: Jossey-Bass,
uma propriedade comum a todos os sistemas, é que o conhecimento de uma característica de
alguma sua parte, estando ela relacionada com outras partes ou com o sistema como um todo,
permite por inferência, algum conhecimento dessas partes ou do sistema;
-For a mathematical dynamical system, S.Smale, proposed the concept of phase space, where
if the system changed, a trajectory could be drawn on paper to represent the changing state of
the system. Phase space contains the complete knowledge of a system. Each point in phase
191
space represents the state of a dynamic system at an instant in time.
– De forma análoga, em estatística,, quando do conhecimento de uma propriedade de
uma amostra (convenientemente escolhida), infere-se o resultado para toda a população;
Alfred Kuhn: “Since Descartes, the "scientific method" had progressed under two related
assumptions. A system could be broken down into its individual components so that each
component could be analyzed as an independent entity, and the components could be added in
a linear fashion to describe the totality of the system. Von Bertalanffy proposed that both
assumptions were wrong (not generally valid). On the contrary, a system is characterized by the
interactions of its components and the nonlinearity of those interactions.”
192
DAS AVALIAÇÕES (DE ALUNOS, PROFESSORES E INSTITUIÇÃO) E DA INFORMAÇÃO SOBRE O
DESEMPENHO PROFISSIONAL DOS EX-ALUNOS, OCORRE O APRIMORAMENTO DO ENSINO
NUM SISTEMA EDUCACIONAL, ETC.;
LEARNING THEORY: Similar ideas are found in learning theories that developed from the
same fundamental concepts, emphasizing how understanding results from knowing concepts
both in part and as a whole. Interdisciplinary perspectives are critical in breaking away from
industrial age models and thinking where history is history and math is math, the arts and
sciences specialized and separate, and where teaching is treated as behaviorist conditioning.
The influential contemporary work of Peter Senge[1] provides detailed discussion of the
commonplace critique of educational systems grounded in conventional assumptions about
learning, including the problems with fragmented knowledge and lack of holistic learning from
the "machine-age thinking" that became a "model of school separated from daily life."
[1]Senge, P., Ed. (2000). Schools That Learn: A Fifth Discipline Fieldbook for Educators, Parents,
and Everyone Who Cares About Education. New York: Doubleday Dell Publishing Group.
-Systems analysis is the study of sets of interacting entities, including computer systems
analysis. This field is closely related to requirements analysis or operations research. It is also
"an explicit formal inquiry carried out to help someone (referred to as the decision maker)
identify a better course of action and make a better decision than he might otherwise have
made.
System development can generally be thought of having two major components: systems
analysis and systems design. In System Analysis more emphasis is given to understanding the
details of an existing system or a proposed one and then deciding whether the proposed
system is desirable or not and whether the existing system needs improvements. Thus, system
analysis is the process of investigating a system, identifying problems, and using the
information to recommend improvements to the system.
-System analysis in the field of electrical engineering characterizes electrical systems and
their properties. System Analysis can be used to represent almost anything from population
growth to audio speakers, electrical engineers often use it because of its direct relevance to
many areas of their discipline, most notably signal processing and communication systems.
193
Characterization of systems
It is often useful (or necessary) to break up a system into smaller pieces for analysis.
Therefore, we can regard a SIMO system as multiple SISO systems (one for each
output), and similarly for a MIMO system. By far, the greatest amount of work in
system analysis has been with SISO systems, although many parts inside SISO systems
have multiple inputs (such as adders).
Signals that are continuous in time and continuous in value are known as analog
signals.
Signals that are discrete in time and discrete in value are known as digital signals.
Signals that are discrete in time and continuous in value are called discrete-time
signals. While important mathematically, systems that process discrete time signals
are difficult to physically realize. The methods developed for analyzing discrete time
signals and systems are usually applied to digital and analog signals and systems.
Signals that are continuous in time and discrete in value are sometimes seen in the
timing analysis of logic circuits, but have little to no use in system analysis.
With this categorization of signals, a system can then be characterized as to which type
of signals it deals with:
A system that has analog input and analog output is known as an analog system.
A system that has digital input and digital output is known as a digital system.
Systems with analog input and digital output or digital input and analog output are
possible. However, it is usually easiest to break these systems up for analysis into their
analog and digital parts, as well as the necessary analog to digital or digital to analog
converter.
Another way to characterize systems is by whether their output at any given time
depends only on the input at that time or perhaps on the input at some time in the past
(or in the future!).
194
Note: It is not possible to physically realize a non-causal system operating in "real
time". However, from the standpoint of analysis, they are important for two reasons.
First, the ideal system for a given application is often a noncausal system, which
although not physically possible can give insight into the design of a derivated causal
system to accomplish a similar purpose. Second, there are instances when a system
does not operate in "real time" but is rather simulated "off-line" by a computer, such
as post-processing an audio or video recording.
Further, some non-causal systems can operate in pseudo-real time by introducing lag:
if a system depends on input for 1 second in future, it can process in real time with 1
second lag.
Analog systems with memory may be further classified as lumped or distributed. The
difference can be explained by considering the meaning of memory in a system. Future
output of a system with memory depends on future input and a number of state
variables, such as values of the input or output at various times in the past. If the
number of state variables necessary to describe future output is finite, the system is
lumped; if it is infinite, the system is distributed.
A system is linear if it has the superposition and scaling properties. A system that is not
linear is non-linear.
If the output of a system does not depend explicitly on time, the system is said to be
time-invariant; otherwise it is time-variant
A system that will always produce the same output for a given input is said to be
deterministic.
A system that will produce different outputs for a given input is said to be non-
deterministic (including stochastic).
There are many methods of analysis developed specifically for linear time-invariant
(LTI) deterministic systems. Unfortunately, in the case of analog systems, none of these
properties are ever perfectly achieved. Linearity implies that operation of a system can
be scaled to arbitrarily large magnitudes, which is not possible. Time-invariance is
violated by aging effects that can change the outputs of analog systems over time
(usually years or even decades). Thermal noise and other random phenomena ensure
that the operation of any analog system will have some degree of stochastic behavior.
Despite these limitations, however, it is usually reasonable to assume that deviations
from these ideals will be small.
LTI Systems
Main article: LTI system theory
As mentioned above, there are many methods of analysis developed specifically for LTI
systems. This is due to their simplicity of specification. An LTI system is completely
specified by its transfer function (which is a rational function for digital and lumped
analog LTI systems). Alternatively, we can think of an LTI system being completely
specified by its frequency response. A third way to specify an LTI system is by its
characteristic linear differential equation (for analog systems) or linear difference
195
equation (for digital systems). Which description is most useful depends on the
application.
The distinction between lumped and distributed LTI systems is important. A lumped
LTI system is specified by a finite number of parameters, be it the zeros and poles of its
transfer function, or the coefficients of its differential equation, whereas specification of
a distributed LTI system requires a complete function
EXEMPLOS:
ONDE
Time delays of one type or another have been incorporated into biological model to represent
resource regeneration times, maturation periods, feeding times, reaction times, etc. by many
researchers.
Delay-differential equations (DDEs) are a large and important class of dynamical systems. They
often arise in either natural or technological control problems. In these systems, a controller
monitors the state of the system, and makes adjustments to the system based on its
observations. Since these adjustments can never be made instantaneously, a delay arises
between the observation and the control action.
Delay differential equations (DDEs) have been used for many years in control theory and only
recently have been applied to biological models. Most biological systems have time delays
inherent in them (for example, epidemiology ), yet few scientists apply these equations due to
the complexity they introduce.
196
para .
Smale, S. 1980. "How I got started in dynamical systems." In The Mathematics of Time:
Essays on Dynamical Systems, Economic Processes, and Related Topics. Smale, Yorke,
Guckenheimer, Abraham, May, Feigenbaum. New York: Springer-Verlag. p. 147-151.,
197
proposed the concept of phase space, where if the system changed, a trajectory could be
drawn on paper to represent the changing state of the system. Phase space contains the
complete knowledge of a system. Each point in phase space represents the state of a dynamic
system at an instant;
A dynamical system is a concept in mathematics where a fixed rule describes the time
dependence of a point in a geometrical space. Examples include the mathematical
models that describe the swinging of a clock pendulum, the flow of water in a pipe, and
the number of fish each springtime in a lake.
At any given time a dynamical system has a state given by a set of real numbers (a
vector) that can be represented by a point in an appropriate state space (a geometrical
manifold). Small changes in the state of the system create small changes in the numbers.
The evolution rule of the dynamical system is a fixed rule that describes what future
states follow from the current state. The rule is deterministic; in other words, for a given
time interval only one future state follows from the current state.
The concept of a dynamical system has its origins in Newtonian mechanics. There, as in
other natural sciences and engineering disciplines, the evolution rule of dynamical
systems is given implicitly by a relation that gives the state of the system only a short
time into the future. (The relation is either a differential equation, difference equation or
other time scale.) To determine the state for all future times requires iterating the
relation many times—each advancing time a small step. The iteration procedure is
referred to as solving the system or integrating the system. Once the system can be
solved, given an initial point it is possible to determine all its future positions, a
collection of points known as a trajectory or orbit.
Before the advent of fast computing machines, solving a dynamical system required
sophisticated mathematical techniques and could be accomplished only for a small class
of dynamical systems. Numerical methods implemented on electronic computing
machines have simplified the task of determining the orbits of a dynamical system.
For simple dynamical systems, knowing the trajectory is often sufficient, but most
dynamical systems are too complicated to be understood in terms of individual
trajectories. The difficulties arise because:
198
the qualitative study of dynamical systems, that is, properties that do not change
under coordinate changes. Linear dynamical systems and systems that have two
numbers describing a state are examples of dynamical systems where the possible
classes of orbits are understood.
The behavior of trajectories as a function of a parameter may be what is needed for an
application. As a parameter is varied, the dynamical systems may have bifurcation
points where the qualitative behavior of the dynamical system changes. For example,
it may go from having only periodic motions to apparently erratic behavior, as in the
transition to turbulence of a fluid.
The trajectories of the system may appear erratic, as if random. In these cases it may
be necessary to compute averages using one very long trajectory or many different
trajectories. The averages are well defined for ergodic systems and a more detailed
understanding has been worked out for hyperbolic systems. Understanding the
probabilistic aspects of dynamical systems has helped establish the foundations of
statistical mechanics and of chaos.
It was in the work of Poincaré that these dynamical systems themes developed.
Generic properties:
From Wikipedia, the free encyclopedia
In mathematics, properties that hold for "typical" examples are called generic
properties. For instance, a generic property of a class of functions is one that is true of
"almost all" of those functions, as in the statements, "A generic polynomial does not
have a root at zero," or "A generic matrix is invertible." As another example, a generic
property of a space is a property that holds at "almost all" points of the space, as in the
statement, "If f : M → N is a smooth function between smooth manifolds, then a generic
point of N is not a critical value of f." (This is by Sard's theorem.)
There are many different notions of "generic" (what is meant by "almost all") in
mathematics, with corresponding dual notions of "almost none" (negligible set); the two
main classes are:
In measure theory, a generic property is one that holds almost everywhere, meaning
"with probability 1", with the dual concept being null set, meaning "with probability
0".
In topology and algebraic geometry, a generic property is one that holds on a dense
open set, or more generally on a residual set, with the dual concept being a nowhere
dense set, or more generally a meagre set
199
The final equation for the linearization of a function at is:
point is:
Linearization makes it possible to use tools for studying linear systems to analyze the
behavior of a nonlinear function near a given point. The linearization of a function is the
first order term of its Taylor expansion around the point of interest. For a system
defined by the equation
at .
Stability analysis
200
In stability analysis, one can use the eigenvalues of the Jacobian matrix evaluated at an
equilibrium point to determine the nature of that equilibrium. If all of the eigenvalues
are positive, the equilibrium is unstable; if they are all negative the equilibrium is stable;
and if the values are of mixed signs, the equilibrium is possibly a saddle point. Any
complex eigenvalues will appear in complex conjugate pairs and indicate a spiral.
Microeconomics
STABILITY
Stability theory, the study of the stability of solutions to differential equations and
dynamical systems
o Lyapunov stability
o Structural stability
Stability (probability), a property of probability distributions
Stability (learning theory), a property of machine learning algorithms
Numerical stability, a property of numerical algorithms which describes how errors in
the input data propagate through the algorithm
Stability radius, a property of continuous polynomial functions
Stable theory, concerned with the notion of stability in model theory
Stability, a property of points in geometric invariant theory
Many parts of the qualitative theory of differential equations and dynamical systems
deal with asymptotic properties of solutions and the trajectories—what happens with the
system after a long period of time. The simplest kind of behavior is exhibited by
equilibrium points, or fixed points, and by periodic orbits. If a particular orbit is well
understood, it is natural to ask next whether a small change in the initial condition will
lead to similar behavior. Stability theory addresses the following questions: will a
nearby orbit indefinitely stay close to a given orbit? will it converge to the given orbit
(this is a stronger property)? In the former case, the orbit is called stable and in the
latter case, asymptotically stable, or attracting. Stability means that the trajectories do
not change too much under small perturbations. The opposite situation, where a nearby
orbit is getting repelled from the given orbit, is also of interest. In general, perturbing
the initial state in some directions results in the trajectory asymptotically approaching
the given one and in other directions to the trajectory getting away from it. There may
also be directions for which the behavior of the perturbed orbit is more complicated
(neither converging nor escaping completely), and then stability theory does not give
sufficient information about the dynamics.
One of the key ideas in stability theory is that the qualitative behavior of an orbit under
perturbations can be analyzed using the linearization of the system near the orbit. In
particular, at each equilibrium of a smooth dynamical system with an n-dimensional
phase space, there is a certain n×n matrix A whose eigenvalues characterize the
201
behavior of the nearby points (Hartman-Grobman theorem). More precisely, if all
eigenvalues are negative real numbers or complex numbers with negative real parts then
the point is a stable attracting fixed point, and the nearby points converge to it at an
exponential rate, cf Lyapunov stability and exponential stability. If none of the
eigenvalues is purely imaginary (or zero) then the attracting and repelling directions are
related to the eigenspaces of the matrix A with eigenvalues whose real part is negative
and, respectively, positive. Analogous statements are known for perturbations of more
complicated orbits.
Various types of stability may be discussed for the solutions of differential equations
describing dynamical systems. The most important type is that concerning the stability of
solutions near to a point of equilibrium. This may be discussed by the theory of Lyapunov. In
simple terms, if all solutions of the dynamical system that start out near an equilibrium point
stay near forever, then is Lyapunov stable. More strongly, if is Lyapunov stable
and all solutions that start out near converge to , then is asymptotically stable.
The notion of exponential stability guarantees a minimal rate of decay, i.e., an estimate of
how quickly the solutions converge. The idea of Lyapunov stability can be extended to infinite-
dimensional manifolds, where it is known as structural stability, which concerns the behavior
of different but "nearby" solutions to differential equations. Input-to-state stability (ISS)
applies Lyapunov notions to systems with inputs.
202
unstable equilibrium, such as a ball resting on a top of a hill, certain small pushes will
result in a motion with a large amplitude that may or may not converge to the original
state.
There are useful tests of stability for the case of a linear system. Stability of a nonlinear
system can often be inferred from the stability of its linearization.
where denotes the system state vector, an open set containing the
origin, and continuous on . Suppose has an equilibrium .
1. The equilibrium of the above system is said to be Lyapunov stable, if, for every
, there exists a such that, if , then
, for every .
2. The equilibrium of the above system is said to be asymptotically stable if it is
Lyapunov stable and if there exists such that if , then
.
3. The equilibrium of the above system is said to be exponentially stable if it is
asymptotically stable and if there exist such that if
, then , for .
203
1. Lyapunov stability of an equilibrium means that solutions starting "close enough" to
the equilibrium (within a distance from it) remain "close enough" forever (within a
distance from it). Note that this must be true for any that one may want to choose.
2. Asymptotic stability means that solutions that start close enough not only remain close
enough but also eventually converge to the equilibrium.
3. Exponential stability means that solutions not only converge, but in fact converge
faster than or at least as fast as a particular known rate .
for for all trajectories that start close enough, and globally attractive if this
property holds for all trajectories.
That is, if x belongs to the interior of its stable manifold. It is asymptotically stable if it
is both attractive and stable. (There are counterexamples showing that attractivity does
not imply asymptotic stability. Such examples are easy to create using homoclinic
connections.)
Lyapunov, in his original 1892 work, proposed two methods for demonstrating stability.
The first method developed the solution in a series which was then proved convergent
within limits. The second method, which is almost universally used nowadays, makes
use of a Lyapunov function V(x) which has an analogy to the potential function of
classical dynamics. It is introduced as follows for a system having a point of
equilibrium at x=0. Consider a function such that
Then V(x) is called a Lyapunov function candidate and the system is asymptotically
stable in the sense of Lyapunov (i.s.L.). (Note that is required; otherwise
for example would "prove" that is locally stable.
An additional condition called "properness" or "radial unboundedness" is required in
order to conclude global asymptotic stability.)
204
Lyapunov's realization was that stability can be proven without requiring knowledge of
the true physical energy, provided a Lyapunov function can be found to satisfy the
above constraints.
The definition for discrete-time systems is almost identical to that for continuous-time
systems. The definition below provides this, using an alternate language commonly
used in more mathematical texts.
then
for all .
We say that is asymptotically stable if it belongs to the interior of its stable set, i.e. if
there is a such that
Whenever .
is asymptotically stable (in fact, exponentially stable) if all real parts of the eigenvalues
205
Correspondingly, a time-discrete linear state space model
is asymptotically stable (in fact, exponentially stable) if all the eigenvalues of have a
modulus smaller than one.
This latter condition has been generalized to switched systems: a linear switched
discrete time system (ruled by a set of matrices )
is asymptotically stable (in fact, exponentially stable) if the joint spectral radius of the
set is smaller than one.
where the (generally time-dependent) input u(t) may be viewed as a control, external
input, stimulus, disturbance, or forcing function. The study of such systems is the
subject of control theory and applied in control engineering. For systems with inputs,
one must quantify the effect of inputs on the stability of the system. The main two
approaches to this analysis are BIBO stability (for linear systems) and input-to-state
(ISS) stability (for nonlinear systems).
206
problem. The leading term in this power series is the solution of the exactly solvable
problem, while further terms describe the deviation in the solution, due to the deviation
from the initial problem. Formally, we have for the approximation to the full solution A,
a series in the small parameter (here called ), like the following:
In this example, would be the known solution to the exactly solvable initial problem
and , ... represent the higher-order terms which may be found iteratively by
some systematic procedure. For small these higher-order terms in the series become
successively smaller. An approximate "perturbation solution" is obtained by truncating
the series, usually by keeping only the first two terms, the initial solution and the "first-
order" perturbation correction:
h) CHAOS THEORY studies the behavior of dynamical systems that are highly sensitive
to initial conditions; For chaotic systems, small differences in initial conditions may
yield widely diverging outcomes, rendering long-term prediction impossible in general.
This happens even though these systems are deterministic, meaning that their future
behavior is fully determined by their initial conditions, with no random elements
involved. In other words, the deterministic nature of these systems does not make them
predictable. Some researchers say that chaotic behavior has an apparently random
outcome.
It can be difficult to tell from data whether a physical or other observed process is
random or chaotic, because in practice no time series consists of pure 'signal.' There will
always be some form of corrupting noise, even if it is present as round-off or truncation
error. Thus any real time series, even if mostly deterministic, will contain some
randomness.
All methods for distinguishing deterministic and stochastic processes rely on the fact
that a deterministic system always evolves in the same way from a given starting point.
Thus, given a time series to test for determinism, one can:
Define the error as the difference between the time evolution of the 'test' state and the
time evolution of the nearby state. A deterministic system will have an error that either
remains small (stable, regular solution) or increases exponentially with time (chaos). A
stochastic system will have a randomly distributed error.
Essentially all measures of determinism taken from time series rely upon finding the
closest states to a given 'test' state (e.g., correlation dimension, Lyapunov exponents,
etc.). To define the state of a system one typically relies on phase space embedding
207
methods. Typically one chooses an embedding dimension, and investigates the
propagation of the error between two nearby states. If the error looks random, one
increases the dimension. If you can increase the dimension to obtain a deterministic
looking error, then you are done. Though it may sound simple it is not really. One
complication is that as the dimension increases the search for a nearby state requires a
lot more computation time and a lot of data (the amount of data required increases
exponentially with embedding dimension) to find a suitably close candidate. If the
embedding dimension (number of measures per state) is chosen too small (less than the
'true' value) deterministic data can appear to be random but in theory there is no
problem choosing the dimension too large – the method will work.
208
larger parameter space, catastrophe theory reveals that such bifurcation points tend
to occur as part of well-defined qualitative geometrical structures.
Elementary catastrophes
Catastrophe theory analyses degenerate critical points of the potential function — points
where not just the first derivative, but one or more higher derivatives of the potential function
are also zero. These are called the germs of the catastrophe geometries. The degeneracy of
these critical points can be unfolded by expanding the potential function as a Taylor series in
small perturbations of the parameters.
When the degenerate points are not merely accidental, but are structurally stable, the
degenerate points exist as organising centres for particular geometric structures of lower
degeneracy, with critical features in the parameter space around them. If the potential
function depends on two or fewer active variables, and four (resp. five) or fewer active
parameters, then there are only seven (resp. eleven) generic structures for these bifurcation
geometries, with corresponding standard forms into which the Taylor series around the
catastrophe germs can be transformed by diffeomorphism (a smooth transformation whose
inverse is also smooth).
i) POPULATION GROWTH
Simple models of population growth include the Malthusian Growth Model (named after the
Reverend Thomas Malthus, who authored An Essay on the Principle of Population, one of the
earliest and most influential books on population and exponetial growth, 1798 ) and the
logistic model (named after Pierre François Verhulst, 1838/45 );
to give
Where . Exponentiating,
209
This equation is called the law of growth, and the quantity in this equation is sometimes known as the
Malthusian parameter.
From a biological point of view the missing feature of Malthus's model is the idea of
carrying capacity. As a population increases in size the environment's ability to support
the population decreases. As the population increases per capita food availability
decreases, waste products may accumulate and birth rates tend to decline while death
rates tend to increase. It seems reasonable to consider a mathematical model which
explicitly incorporates the idea of carrying capacity.
Note that this expression blows up at t=0 . We are given the initial condition that
N (t=1) = N0 er , so C = N0 , and.
The in the denominator greatly suppresses the growth in the long run compared to the simple growth law.
210
The Logistic Population Model
The logistic model, a slight modification of Malthus's model, is just such a model. As with
Malthus's model the logistic model includes a growth rate r. This parameter represents the
rate at which the population would grow if it were unencumbered by environmental
degradation.
A second parameter, K, represents the carrying capacity of the system being studied.
Carrying capacity is the population level at which the birth and death rates of a species
precisely match, resulting in a stable population over time.
Descrete model: letting X(i) represent the population at the beginning of time period i
the logistic model is:
The first term of this formula, r*X(i), can be interpreted as the birth rate in this model,
while the second term, r*X(i)*X(i)/K can be interpreted as the death rate. Noting that the
birth rate only depends on r while the death rate depends on both r and K, the
appropriate connectors can be created and erased in order to place the required formula
in each flow.
This approach of explicitly determining the birth and death rates in a population model
is routinely implemented because it gives the modeler a chance to look at the model
components to see if the birth and death formulas are reasonable. It also provides
experimentalists with a way to break the larger problem of population dynamics into
two simpler pieces. Instead of having to derive one experiment to measure population
changes as a function of population size, an experimentalist (or theorist for that matter)
can consider the birth process and death process separately.
The analysis which can be performed by separating out the birth and death processes is
well illustrated by the logistic model. The birth rate seems quite reasonable. It depends
only on the per capita birth rate r. It may strike some as odd that the per capita birth rate
does not decline as population increases. This is a weakness in the logistic model, and
various attempts at correcting this problem have given rise to many variations on the
logistic model.
A more serious weakness in the logistic model becomes clear when the death rate
formula is considered. The number of deaths, given by
r*X(i)*X(i)/K,
does increase as the population increases. This represents the effects of congestion. The
carrying capacity K plays a significant role in this formula, as it should. However, it
seems odd that r, the per capita birth rate appears in this formula. It does not seem
reasonable that the birth rate should directly affect the number of deaths experienced in
a population. This is regarded by many as an indication that the logistic model is
probably not a good representation of biological reality.
211
Despite the weaknesses illustrated above, the logistic model is frequently used in
biological modeling, either a basis for more complicated modeling, or as an
approximate model when the details of a population dynamics are not known.
is another growth law which frequently arises in biology. Solving the equation by separating the variables,
one obtains
Hutchinson (1948) assumed egg formation to occur 0 units of time before hatching
and proposed the following more realistic logistic equation
dN/dt = r N(t) { 1 - [ N ( t - n τ ) ] / K }
The Hutchinson's equation means that the regulatory efect depends on the population at a fixed earlier
time t - , rather than at the present time t. In a more realistic model the delay efect should be an average
over past populations. This results in an equation with a distributed delay or an infnite delay. The first
work using a logistic equation with distributed delay was by Volterra (1934) with extensions by Kostitzin
(1939). In the 1930's, many experiments were performed with laboratory populations of some species of
small organisms with short generation time. Attempts to apply logistic models to these experiments were
often unsuccessful because populations died out. One of the causes was the pollution of the closed
environment by waste products and dead organisms. Volterra (1934) used an integral term or a distributed
delay term to examine a cumulative efect in the death rate of a species, depending on the population at all
times from the start of the experiment. The model is an integro-diferential equation:
dN/dt = r N [ 1/k ]
This
is an integro-diferential equation where G(t), called the delay kernel, is a weighting
factor which indicates how much emphasis should be given to the size of the population
at earlier times to determine the present efect on resource availability.
212
A generalized logistic function can model the "S-shaped" behaviour (abbreviated S-
curve) of growth of some population P. The initial stage of growth is approximately
exponential; then, as saturation begins, the growth slows, and at maturity, growth stops.
A simple logistic function may be defined by the formula
OBSERVAÇÃO: O PONTO DE VISTA DOS SISTEMAS PERMITE CONSIDERAR O EFEITO DAS TAXAS
DE MORTALIDADE E MIGRAÇÕES, SOBRE A TAXA DE NASCIMENTOS;
Population[1]
Year Billion
1800 1
1927 2
213
1960 3
1974 4
1987 5
1999 6
2011* 7
∆P≡(B-D)+(I-E)
In demographics and ecology, population growth rate (PGR) is the rate at which the
number of individuals in a population increases in a given time period as a fraction of
the initial population. Specifically, PGR ordinarily refers to the change in population
over a unit time period, often expressed as a percentage of the number of individuals in
the population at the beginning of that period. This can be written as the formula:[2]
The most common way to express population growth is as a percentage, not as a rate.
The change in population over a unit time period is expressed as a percentage of the
population at the beginning of the time period. That is:
If the length of the time is taken smaller and smaller, the PGR approaches the
logarithmic derivative of the population function P. If the population as a function of
time is exponential, say P(t) = Ceat, the logarithmic derivative is a. Thus, the PGR
approximates the exponent a for populations with exponential growth.
Globally, the growth rate of the human population has been declining since peaking in
1962 and 1963 at 2.20% per annum. In 2009, the estimated annual growth rate was
1.1%. The CIA World Factbook gives the world annual birthrate, mortality rate, and
growth rate as 1.915%, 0.812%, and 1.092% respectively. The last one hundred years
have seen a rapid increase in population due to medical advances and massive increase
in agricultural productivity made possible by the Green Revolution.
The actual annual growth in the number of humans fell from its peak of 88.0 million in
1989, to a low of 73.9 million in 2003, after which it rose again to 75.2 million in 2006.
214
Since then, annual growth has declined. In 2009, the human population increased by
74.6 million, which is projected to fall steadily to about 41 million per annum in 2050,
at which time the population will have increased to about 9.2 billion. [11] Each region of
the globe has seen great reductions in growth rate in recent decades, though growth
rates remain above 2% in some countries of the Middle East and Sub-Saharan Africa,
and also in South Asia, Southeast Asia, and Latin America.[12]
Um modelo clássico onde se pode utilizar esta fórmula é aquele que envolve as relações
entre o lince e lebre, ou lobo e coelho;
215
A forma usual da equação é :
onde
The Lotka-Volterra model makes a number of assumptions about the environment and
evolution of the predator and prey populations:
As differential equations are used, the solution is deterministic and continuous. This, in
turn, implies that the generations of both the predator and prey are continually
overlapping.[20]
Prey
The prey are assumed to have an unlimited food supply, and to reproduce exponentially
unless subject to predation; this exponential growth is represented in the equation above
by the term αx. The rate of predation upon the prey is assumed to be proportional
to the rate at which the predators and the prey meet; this is represented above by
βxy. If either x or y is zero then there can be no predation.
216
With these two terms the equation above can be interpreted as: the change in the prey's
numbers is given by its own growth minus the rate at which it is preyed upon.
Predators
In this equation, δxy represents the growth of the predator population. (Note the
similarity to the predation rate; however, a different constant is used as the rate at which
the predator population grows is not necessarily equal to the rate at which it consumes
the prey). γy represents the loss rate of the predators due to either natural death or
emigration; it leads to an exponential decay in the absence of prey.
Hence the equation expresses the change in the predator population as growth fueled by
the food supply, minus natural death.
As populações de presas e
predadores em função do tempo
obtida utilizando o modelo de
Lotka-Volterra. Os parâmetros
utilizados são
The equations have periodic solutions and do not have a simple expression in terms of
the usual trigonometric functions. However, a linearization of the equations yields a
solution similar to simple harmonic motion[21] with the population of predators
following that of prey by 90°.
217
A SOLUÇÃO DO SISTEMA ACIMA PODE SER INTERPRETADA DO SEGUINTE MODO:
The Lotka–Volterra predator–prey model was initially proposed by Alfred J. Lotka “in
the theory of autocatalytic chemical reactions” in 1910. This was effectively the logistic
equation, which was originally derived by Pierre François Verhulst. In 1920 Lotka
extended, via Kolmogorov (see above), the model to "organic systems" using a plant
species and a herbivorous animal species as an example and in 1925 he utilised the
equations to analyse predator-prey interactions in his book on biomathematics arriving
at the equations that we know today. Vito Volterra, who made a statistical analysis of
fish catches in the Adriatic independently investigated the equations in 1926.
C.S. Holling extended this model yet again, in two 1959 papers, in which he proposed
the idea of functional response. Both the Lotka-Volterra model and Holling's
extensions have been used to model the moose and wolf populations in Isle Royale
National Park, which with over 50 published papers is one of the best studied predator-
prey relationships.
In economics
The Lotka–Volterra equations have a long history of use in economic theory; their
initial application is commonly credited to Richard Goodwin in 1965 or 1967. In
economics, links are between many if not all industries; a proposed way to model the
dynamics of various industries has been by introducing trophic functions between
various sectors, and ignoring smaller sectors by considering the interactions of only two
industrial sectors.
218
k) GENE PROPAGATION: Fisher’s equation (or gene transport equation.) appeared in his
seminal paper R. A. Fisher, “The wave of advance of advantageous genes”, Ann. Eugenics
7 (1937) 335–369:
+ m p (1 – p )
Alan Lloyd Hodgkin and Andrew Huxley described the model in 1952 to explain the ionic
mechanisms underlying the initiation and propagation of action potentials in the squid giant
axon. Hodgkin, A., and Huxley, A. (1952): A quantitative description of membrane current and
its application to conduction and excitation in nerve. J. Physiol. 117:500–544
219
They received the 1963 Nobel Prize in Physiology or Medicine for this work.
Huxley equation is a nonlinear partial diferential equation of second order of the form
ut = uxx + u ( - u) (u – 1)
This is an evolution equation that describes the nerve propagation in biology from
which molecular CB properties can be calculated. It also givees a phenomenological
description of the behavior of the myosin heads II. This equation has many fascinating
phenomena such a bursting oscilation, interspike, bifurcation and chaos. There are
many methods to solve the equation, but each of them can only lead to a special
solution. (Mathamatical Problems in Engineering 2008, Exp-Function Method for
Solving Huxley Equation, Xin-Wei Zhou ).
l) Pendulum
The mathematics of pendulums are in general very complex. Simplifying assumptions
can be made, which in the case of a simple pendulum allows the equations of motion
to be solved analytically for small-angle oscillations.
The rod or cord on which the bob swings is massless, inextensible and always
remains taut;
220
Motion occurs only in two dimensions, i.e. the bob does not trace an ellipse but
an arc.
The motion does not lose energy to friction or air resistance.
The differential equation which represents the motion of a simple pendulum is
(Eq1)
where is acceleration due to gravity, is the length of the pendulum, and is the
angular displacement.
"Force" derivation of (Eq. 1)
One can note that the path of the pendulum sweeps out an arc of a circle. The
angle is measured in radians, and this is crucial for this formula. The blue arrow is
the gravitational force acting on the bob, and the violet arrows are that same force
resolved into components parallel and perpendicular to the bob's instantaneous
motion. The direction of the bob's instantaneous velocity always points along the red
axis, which is considered the tangential axis because its direction is always tangent to
the circle. Consider Newton's second law,
where : F is the sum of forces on the object, m is mass, and a is the acceleration.
Because we are only concerned with changes in speed, and because the bob is forced
to stay in a circular path, we apply Newton's equation to the tangential axis only. The
221
short violet arrow represents the component of the gravitational force in the tangential
axis, and trigonometry can be used to determine its magnitude. Thus,
Where: is the acceleration due to gravity near the surface of the earth. The negative
sign on the right hand side implies that and always point in opposite directions. This
makes sense because when a pendulum swings further to the left, we would expect it
to accelerate back toward the right.
This linear acceleration along the red axis can be related to the change in angle by
the arc length formulas; is arc length:
thus:
First start by defining the torque on the pendulum bob using the force due to gravity.
where is the length vector of the pendulum and is the force due to gravity.
For now just consider the magnitude of the torque on the pendulum.
where is the mass of the pendulum, is the acceleration due to gravity, is the
length of the pendulum and is the angle between the length vector and the force due
to gravity.
222
.
thus:
It can also be obtained via the conservation of mechanical energy principle: any object
falling a vertical distance would acquire kinetic energy equal to that which it lost to
the fall. In other words, gravitational potential energy is converted into kinetic energy.
Change in potential energy is given by
223
Using the arc length formula above, this equation can be rewritten in favor of
is the vertical distance the pendulum fell. Look at Figure 2, which presents the
trigonometry of a simple pendulum. If the pendulum starts its swing from some initial
angle , then , the vertical distance from the screw, is given by
in terms of gives
(Eq. 2)
This equation is known as the first integral of motion, it gives the velocity in terms of the
location and includes an integration constant related to the initial displacement ( ).
We can differentiate, by applying the chain rule, with respect to time to get the
acceleration
224
Small-angle approximation
The differential equation given above is not easily solved, and there is no solution that
can be written in terms of elementary functions. However adding a restriction to the
size of the oscillation's amplitude gives a form whose solution can be easily obtained. If
it is assumed that the angle is much less than 1 radian, or
,
then substituting for sin θ into Eq. 1 using the small-angle approximation,
,
yields the equation for a harmonic oscillator,
The error due to the approximation is of order θ 3 (from the Maclaurin series for sin θ).
Given the initial conditions θ(0) = θ0 and dθ/dt(0) = 0, the solution becomes,
The motion is simple harmonic motion where θ0 is the semi-amplitude of the oscillation
(that is, the maximum angle between the rod of the pendulum and the vertical). The
period of the motion, the time for a complete oscillation (outward and return) is
which is known as Christian Huygens's law for the period. Note that under the small-
angle approximation, the period is independent of the amplitude θ0; this is the property
of isochronism that Galileo discovered.
can be expressed as
If SI units are used (i.e. measure in metres and seconds), and assuming the
measurement is taking place on the Earth's surface, then m/s2,
and (0.994 is the approximation to 3 decimal places).
Therefore a relatively reasonable approximation for the length and period are,
225
Arbitrary-amplitude period
For amplitudes beyond the small angle approximation, one can compute the exact
period by first inverting the equation for the angular velocity obtained from the energy
method (Eq. 2),
and then integrating over one complete cycle,
which leads to
(Eq. 3)
For comparison of the approximation to the full solution, consider the period of a pendulum of
length 1 m on Earth (g = 9.80665 m/s2) at initial angle 10 degrees
is .
The difference between the two values, less than 0.2%, is much less than that caused by
the variation of g with geographical location.
226
Legendre polynomial solution for the elliptic integral
Given Eq. 3 and the Legendre polynomial solution for the elliptic integral:
where n!! denotes the double factorial, an exact solution to the period of a
pedulum is:
Figure 4 shows the relative errors using the power series. T0 is the linear
approximation, and T2 to T10 include respectively the terms up to the 2nd to the
10th powers.
227
Power series solution for the elliptic integral
Another formulation of the above solution can be found if the following
Maclaurin series:
is used in the Legendre polynomial solution above. The resulting power series
is:[1]
Given Eq. 3 and the Arithmetic-geometric mean solution of the elliptic integral:
Figure 5. Potential energy and phase portrait of a simple pendulum. Note that the x-
axis, being angle, wraps onto itself after every 2π radians.
228
References
1.^ Nelson, Robert; M. G. Olsson (February 1986). "The pendulum — Rich physics
from a simple system". American Journal of Physics 54 (2): pp. 112–
121. doi:10.1119/1.14703. Retrieved 2012-04-30.
2.^ Carvalhaes, Claudio; Suppes, Patrick (2008). "Approximations for the period of the
simple pendulum based on the arithmetic-geometric mean". American Journal of
Physics 76 (12).
4.^ Paul Appell, "Sur une interprétation des valeurs imaginaires du temps en
Mécanique", Comptes Rendus Hebdomadaires des Scéances de l'Académie des Sciences,
volume 87, number 1, July, 1878
onde e são duas funções quaisquer, que dependem das variáveis e . Não
existem técnicas analíticas gerais para resolver esse tipo de equações; unicamente
existem técnicas analíticas gerais para o caso dos sistemas lineares, em
que e são combinações lineares das variáveis e .1
Os sistemas não lineares geralmente só podem ser resolvidos por métodos numéricos.
No entanto, a análise gráfica no espaço de fase pode fornecer muita informação sobre
o comportamento do sistema.1 Em geral, começa-se por identificar os pontos fixos.
Pontos de equilíbrio
Consideremos o sistema 1:
Existem quatro pontos de equilíbrio. Os pontos onde o lado direito da primeira
Os pontos de equilíbrio são os pontos de interseção entre as curvas onde cada uma
das funções é nula.
229
Os pontos de equilíbrio do sistema são os quatro pontos de interseção entre a elipse e
a hipérbole. Os gráficos dessas duas curvas desenham-se mais facilmente usando a
forma paramétrica dessas equações:
Mudando a origem de coordenadas para o ponto fixo , isto é, num novo sistema
de coordenadas: , , as funções são, aproximadamente,
230
Sistemas Não Lineares Físicos
Pêndulo
O tipo de pêndulo a seguir está formado por um disco de massa e raio , ligado a
uma barra rígida de massa desprezável em comparação com . No outro extremo da
barra passa um eixo horizontal que permite que o pêndulo rode num plano vertical,
descrevendo trajetórias circulares com raio , onde é a distância desde o centro do
disco até o eixo de rotação. (figura abaixo).
Pêndulo formado por um disco e uma barra que pode rodar à volta de um eixo
horizontal.
O pêndulo tem unicamente um grau de liberdade, que pode ser definido como o
ângulo que faz com a vertical. Portanto, existem duas variáveis de estado, , e a
velocidade angular . A primeira equação de evolução é a relação entre o ângulo e a
velocidade angular: . A segunda equação de evolução é a expressão da
aceleração angular em função de e de . Para encontrar essa expressão, é
preciso resolver as leis do movimento do corpo rígido.1
Sobre o pêndulo atuam duas forças externas: o peso, , vertical, e uma força de
contato do eixo sobre a barra, , que por conveniência será decomposta numa
componente tangencial e outra componente normal , na direção da barra.1
Como o eixo de rotação do pêndulo está fixo, pode aplicar-se a lei do movimento de
rotação com eixo fixo:
Neste caso, o peso é a única força que produz momento em relação ao eixo e esse
momento é . Assim, a expressão para em função do ângulo é:
231
e usando o teorema dos eixos paralelos para deslocar o eixo
uma distância , desde o centro do disco até o eixo do pêndulo, obtemos:
O chamado pêndulo simples corresponde ao caso em que o raio do disco, , for muito
menor que o comprimento da barra, ; nesse caso, o momento de inércia será,
aproximadamente, , e as equações de evolução obtidas para o pêndulo
simples são as seguintes (num pêndulo que não seja simples, deverá ser
substituída por
232
Perto do ponto de equilíbrio em , a matriz jacobiana é igual a:
Referência
1. [ Introdução aos sistemas dinâmicos. Porto: Jaime E. Villate, 27 de fevereiro de 2007. 204 págs].
2.^ Carvalhaes, Claudio; Suppes, Patrick (2008). "Approximations for the period of the
simple pendulum based on the arithmetic-geometric mean". American Journal of
Physics 76 (12).
4.^ Paul Appell, "Sur une interprétation des valeurs imaginaires du temps en
Mécanique", Comptes Rendus Hebdomadaires des Scéances de l'Académie des Sciences,
volume 87, number 1, July, 1878
234
COMPOUND PENDULUM
A compound pendulum (or physical pendulum) is one where the rod is not
massless, and may have extended size; that is, an arbitrarily shaped rigid
body swinging by a pivot. In this case the pendulum's period depends on
its moment of inertia I around the pivot point.
The equation of torque gives:
where:
is the angular acceleration.
is the torque
The torque is generated by gravity so:
where:
L is the distance from the pivot to the center of mass of the pendulum
θ is the angle from the vertical
Hence, under the small-angle approximation ,
This is of the same form as the conventional simple pendulum and this gives a
period of:
[3]
235
Double pendulum
A double pendulum is a pendulum with another pendulum attached to its end, and is
a simple physical system that exhibits rich dynamic behavior with a strong sensitivity to
initial conditions.[1] The motion of a double pendulum is governed by a set of
coupled ordinary differential equations. For certain energies its motion is chaotic.
236
and the center of mass of the second pendulum is at
Lagrangian
The Lagrangian is
The first term is the linear kinetic energy of the center of mass of the bodies and the
second term is the rotational kinetic energy around the center of mass of each rod. The
last term is the potential energy of the bodies in a uniform gravitational field. The dot-
notation indicates the time derivative of the variable in question.
Substituting the coordinates above and rearranging the equation gives
There is only one conserved quantity (the energy), and no conserved momenta. The
two momenta may be written as
and
and
237
and
These last four equations are explicit formulae for the time evolution of the system
given its current state. It is not possible to go further and integrate these equations
analytically, to get formulae for θ1 and θ2 as functions of time. It is however possible to
perform this integration numerically using the Runge Kutta method or similar
techniques.
Chaotic motion
The double pendulum undergoes chaotic motion, and shows a sensitive dependence
on initial conditions. The image to the right shows the amount of elapsed time before
the pendulum "flips over," as a function of initial conditions. Here, the initial value of
θ1 ranges along the x-direction, from −3 to 3. The initial value θ2 ranges along the y-
direction, from −3 to 3. The colour of each pixel indicates whether either pendulum flips
then it is energetically impossible for either pendulum to flip. Outside this region, the
pendulum can flip, but it is a complex question to determine when it will flip.
The lack of a natural excitation frequency has led to the use of double pendulum
systems in seismic resistance designs in buildings, where the building itself is the
primary inverted pendulum, and a secondary mass is connected to complete the
double pendulum.
238
Motion of the double compound pendulum (from numerical integration of the equations
of motion) exhibiting chaotic motion (tracked with an LED).
Notes
Meirovitch, Leonard (1986). Elements of Vibration Analysis (2nd edition ed.). McGraw-Hill
Science/Engineering/Math. ISBN 0-07-041342-8.
Eric W. Weisstein, Double pendulum (2005), Science World (contains details of the complicated
equations involved) and "Double Pendulum" by Rob Morris, Wolfram Demonstrations Project, 2007
(animations of those equations).
Peter Lynch, Double Pendulum, (2001). (Java applet simulation.)
Northwestern University, Double Pendulum, (Java applet simulation.)
Theoretical High-Energy Astrophysics Group at UBC, Double pendulum, (2005).
External links
Animations and explanations of a double pendulum and a physical double pendulum (two square
plates) by Mike Wheatland (Univ. Sydney)
Video of a double square pendulum with three (almost) identical starting conditions.
Double pendulum physics simulation from www.myphysicslab.com
239
m) LORENZ SYSTEM
Overview
In 1963, Edward Lorenz developed a simplified mathematical model for atmospheric
convection.[1]
The model is a system of three ordinary differential equations now known as the
Lorenz equations:
Here , , and make up the system state, is time, and , , are the system
parameters. The Lorenz equations also arise in simplified models
[2] [3] [4] [5]
for lasers, dynamos, thermosyphons, brushless DC motors, electric
circuits,[6] and chemical reactions.[7]
From a technical standpoint, the Lorenz system is nonlinear, three-dimensional
and deterministic. The Lorenz equations have been the subject of at least one book
length study.[8]
240
A plot of a solution when ρ = 28,
σ = 10, and β = 8/3 (i.e. a
solution in the Lorenz attractor)
A trajectory of Lorenz's
equations, rendered as a
metal wire to show
direction and 3D structure
Analysis
One normally assumes that the parameters , , and are positive. Lorenz used the
values , and . The system exhibits chaotic behavior for these
values.[9]
If then there is only one equilibrium point, which is at the origin. This point corresponds
to no convection. All orbits converge to the origin when .[10]
A saddle-node bifurcation occurs at , and for two additional critical points appear
at
241
These correspond to steady convection. This pair of equilibrium points is stable only if
which can hold only for positive if . At the critical value, both equilibrium points
lose stability through a Hopf bifurcation.[11]
When , , and , the Lorenz system has chaotic solutions (but not
all solutions are chaotic). The set of chaotic solutions make up the Lorenz attractor, a strange
attractor and a fractal with a Hausdorff dimension which is estimated to be 2.06 ± 0.01 and the
The Lorenz attractor is difficult to analyze, but the action of the differential equation on the
attractor is described by a fairly simple geometric model, proving that this is indeed the case is
the fourteenth problem on the list of Smale's problems. This problem was the first one to be
resolved, by Warwick Tucker in 2002.[13]
For other values of , the system displays knotted periodic orbits. For example,
with it becomes a T(3,2) torus knot.
242
ρ=15, σ=10, β=8/3 ρ=28, σ=10, β=8/3
For small values of ρ, the system is stable and evolves to one of two fixed point attractors. When ρ is
larger than 24.74, the fixed points become repulsors and the trajectory is repelled by them in a very
complex way.
These figures — made using ρ=28, σ = 10 and β = 8/3 — show three time segments of the 3-D evolution
of 2 trajectories (one in blue, the other in yellow) in the Lorenz attractor starting at two initial points that
differ only by 10-5 in the x-coordinate. Initially, the two trajectories seem coincident (only the yellow one
can be seen, as it is drawn over the blue one) but, after some time, the divergence is obvious.
Notes
1. ^ Lorenz (1963)
2. ^ Haken (1975)
243
3. ^ Knobloch (1981)
4. ^ Gorman, Widmann & Robbins (1986)
5. ^ Hemati (1994)
6. ^ Cuomo & Oppenheim (1993)
7. ^ Poland (1993)
8. ^ Sparrow (1982)
9. ^ Hirsch, Smale & Devaney (2003), pp. 303–305
10. ^ Hirsch, Smale & Devaney (2003), pp. 306+307
11. ^ Hirsch, Smale & Devaney (2003), pp. 307+308
12. ^ Grassberger & Procaccia (1983)
13. ^ Tucker (2002)
14. ^ Hilborn (2000), Appendix C; Bergé, Pomeau & Vidal (1984), Appendix D
15. ^ Saltzman (1962)
References
Bergé, Pierre; Pomeau, Yves; Vidal, Christian (1984). Order within Chaos: Towards a
Deterministic Approach to Turbulence. New York: John Wiley & Sons. ISBN 978-0-471-84967-4.
Cuomo, Kevin M.; Oppenheim, Alan V. (1993). "Circuit implementation of synchronized chaos
with applications to communications". Physical Review Letters 71 (1): 65–
68.doi:10.1103/PhysRevLett.71.65. ISSN 0031-9007.
Gorman, M.; Widmann, P.J.; Robbins, K.A. (1986). "Nonlinear dynamics of a convection loop: A
quantitative comparison of experiment with theory". Physica D 19 (2): 255–267.doi:10.1016/0167-
2789(86)90022-9.
Grassberger, P.; Procaccia, I. (1983). "Measuring the strangeness of strange attractors". Physica
D 9 (1–2): 189–208. Bibcode:1983PhyD....9..189G. doi:10.1016/0167-2789(83)90298-1.
Haken, H. (1975). "Analogy between higher instabilities in fluids and lasers". Physics Letters
A 53 (1): 77–78. doi:10.1016/0375-9601(75)90353-9.
Hemati, N. (1994). "Strange attractors in brushless DC motors". IEEE Transactions on Circuits
and Systems I: Fundamental Theory and Applications 41 (1): 40–
45.doi:10.1109/81.260218. ISSN 1057-7122.
Hilborn, Robert C. (2000). Chaos and Nonlinear Dynamics: An Introduction for Scientists and
Engineers (second ed.). Oxford University Press. ISBN 978-0-19-850723-9.
Hirsch, Morris W.; Smale, Stephen; Devaney, Robert (2003). Differential Equations, Dynamical
Systems, & An Introduction to Chaos (Second ed.). Boston, MA: Academic Press. ISBN 978-0-12-
349703-1.
Lorenz, Edward Norton (1963). "Deterministic nonperiodic flow". Journal of the Atmospheric
Sciences 20 (2): 130–141. Bibcode:1963JAtS...20..130L. doi:10.1175/1520-
0469(1963)020<0130:DNF>2.0.CO;2.
Knobloch, Edgar (1981). "Chaos in the segmented disc dynamo". Physics Letters A 82 (9): 439–
440. doi:10.1016/0375-9601(81)90274-7.
Poland, Douglas (1993). "Cooperative catalysis and chemical chaos: a chemical model for the
Lorenz equations". Physica D 65 (1): 86–99. doi:10.1016/0167-2789(93)90006-M.
Saltzman, Barry (1962). "Finite Amplitude Free Convection as an Initial Value Problem—
I". Journal of the Atmospheric Sciences 19 (4): 329–341.
Sparrow, Colin (1982). The Lorenz Equations: Bifurcations, Chaos, and Strange Attractors.
Springer.
Tucker, Warwick (2002). "A Rigorous ODE Solver and Smale's 14th Problem". Foundations of
Computational Mathematics 2 (1): 53–117. doi:10.1007/s002080010018.
244
This is a ranking list of chaotic maps. Only top dozen Best Chaotic Maps, according to
different criteria, are presented. The ranking is in particular based on the number of occurences
of each chaotic map in web pages, news, pictures and people votes in corresponding context.
2. 79 Cellular automata
3. 75 Mandelbrot set
4. 60 Standard map
5. 57 Lorenz attractor
6. 39 Area-preserving maps
8. 38 Quadratic map
9. 37 Cantor set
245
78- CÁLCULO DE PROBABILIDADES. QUANDO HÁ VARIÁVEIS QUE NÃO SÃO, NÃO
CONSEGUEM OU NÃO PODEM SER CONTROLADAS, OS EXPERIMENTOS SÃO MENOS
CONCLUSIVOS, DEPENDENDO DE VÁRIAS REPETIÇÕES, SENDO MUITAS VEZES UTILIZADAS AS
CONCLUSÕES PROBABILÍSTICAS, QUE PERMITEM, AS CORRESPONDENTES PREDIÇÕES
PROBABILÍSTICAS, E CONFORME O CASO, AS ANTECIPAÇÕES, ADAPTAÇÕES, OU AJUSTES,
CONVENIENTES.
, através de:
246
OU SEJA, É A MÉDIA DOS VALORES, PONDERADA PELAS RESPECTIVAS PROBABILIDADES DE
OCORRÊNCIA;
.= , onde μ = E(X)
68% dos valores encontram-se a uma distância da média inferior a um desvio padrão.
95% dos valores encontram-se a uma distância da média inferior a duas vezes o desvio
padrão.
99,7% dos valores encontram-se a uma distância da média inferior a três vezes o
desvio padrão.
247
COLESTEROL, PRESSÃO SANGUÍNEA, FUMO, EXCESSO DE PESO, EXCESSO DE AÇÚCAR NO
SANGUE, VIDA SEDENTÁRIA, ALÉM DE POSSÍVEIS OUTRAS COMO CONSUMO DE FIBRAS, DE
BEBIDAS ALCOÓLICAS, ETC.), QUE, UMA VEZ CONTROLADAS, DIMINUEM O RISCO DE SOFRER
ATAQUES; CONSTATOU-SE AINDA, QUE HÁ PESSOAS QUE PARECERAM IMUNES AOS FATORES
DE RISCO, O QUE SÓ PODE SER CONSTATADO DEPOIS DOS 80 ANOS DE IDADE, QUANDO O
INDIVÍDUO JÁ TERIA CORRIDO O RISCO DE ATAQUE; OBS: A QUANTIDADE E INTENSIDADE DOS
EXERCÍCIOS SÃO AINDA PROBLEMAS ABERTOS (HARVARD ALUMNI HEALTH STUDY);
.
d) Num exemplo divulgado, diz-se que uma consultoria contratada por um empresário,
na pesquisa de dados sobre o mercado, para reconhecer produtos nos quais pudesse
investir, constatou, inesperadamente, que a melhor correlação numérica
(aproximadamente igual a 1), se verificou entre o aumento de salário dos professores
universitários e o consumo de uísque, induzindo a um investindo na marca
Teacher's.
e) EXPECTATIVA DE VIDA
248
The most commonly used measure of life expectancy is life expectancy at age
zero, that is, at birth (LEB), which can be defined as the mean length of life (or
the average number of years) a hypothetical birth cohort (all individuals born a
given year) in a given country would live if mortality rates at each age were to
remain constant in the future; LEB can be computed only for cohorts that were
born many decades ago, so that all their members died. It was estimated that in
the Bronze and Iron Age LEB was 26 years; the 2010 world LEB was 67.2, and
for recent years in Swaziland LEB is about 49 years while in Japan is about 83
years. Worldwide, the average life expectancy at birth was 71.0 years (68.5
years for males and 73.5 years for females) over the period 2010–2013
according to United Nations World Population Prospects 2012 Revision. Os
dados de 2012 da Organização mundial da Saúde da ONU, confirmam que, em
cada país, a expectativa de vida das mulheres e maior do que a dos homens,
ou seja, elas vivem mais e ainda reclamam! O Brasil comparece na 58ª
posição do ranking mundial. The combination of high infant mortality and deaths
in young adulthood from accidents, epidemics, plagues, wars, and childbirth,
particularly before modern medicine was widely available, significantly lowers
LEB. Segundo o IBGE (Instituto Brasileiro de Geografia e Estatística) a
expectativa de vida ao nascer no Brasil (para todas as idades até 80 anos)
subiu de 74,6 anos em 2012, para 74,9 anos em 2013, apresentando um
aumento de 3 meses e 25 dias. Mas, se comparada com a de dez anos atrás, a
expectativa de vida do brasileiro aumentou mais de três anos. Em 2003, era de
71,3 anos.
249
países desenvolvidos, a expectativa de vida (em média), em 1970 era de 68
anos, e hoje já alcança mais de 82 anos. Tentar identificar as alternativas
verdadeiras:
250
Instead of describing a process which can only evolve in one way (as in the case, for example,
of solutions of an ordinary differential equation), in a stochastic or random process there is
some indeterminacy: even if the initial condition (or starting point) is known, there are several
(often infinitely many) directions in which the process may evolve.
In the simple case of discrete time, a stochastic process amounts to a sequence of random
variables known as a time series (for example, see Markov chain). Another basic type of a
stochastic process is a random field, whose domain is a region of space, in other words, a
random function whose arguments are drawn from a range of continuously changing values.
One approach to stochastic processes treats them as functions of one or several deterministic
arguments (inputs, in most cases regarded as time) whose values (outputs) are random variables:
non-deterministic (single) quantities which have certain probability distributions. Random
variables corresponding to various times (or points, in the case of random fields) may be
completely different. The main requirement is that these different random quantities all have the
same type. Type refers to the codomain of the function. Although the random values of a
stochastic process at different times may be independent random variables, in most commonly
considered situations they exhibit complicated statistical correlations.
Familiar examples of processes modeled as stochastic time series include stock market and
exchange rate fluctuations, signals such as speech, audio and video, medical data such as a
patient's EKG, EEG, blood pressure or temperature, and random movement such as Brownian
motion or random walks. Examples of random fields include static images, random terrain
(landscapes), wind waves or composition variations of a heterogeneous material.
A MARKOV CHAIN, named after Andrey Markov, is a mathematical system that undergoes
transitions from one state to another, between a finite or countable number of possible states. It
is a random process characterized as memoryless: the next state depends only on the current
state and not on the sequence of events that preceded it. This specific kind of "memorylessness"
is called the Markov property. Markov chains have many applications as statistical models of
real-world processes.
Formal definition: A Markov chain is a sequence of random variables X1, X2, X3, ... with the
Markov property, namely that, given the present state, the future and past states are independent.
Formally,
251
The possible values of Xi form a countable set S called the state space of the chain.
Markov chains are often described by a directed graph, where the edges are labeled by the
probabilities of going from one state to the other states.
Example:
A simple example is shown in the figure on the right, using a directed graph to picture the state
transitions. The states represent whether the economy is in a bull market, a bear market, or a
recession, during a given week. According to the figure, a bull week is followed by another bull
week 90% of the time, a bear market 7.5% of the time, and a recession the other 2.5%. From
this figure it is possible to calculate, for example, the long-term fraction of time during which
the economy is in a recession, or on average how long it will take to go from a recession to a
bull market. Using the transition probabilities, the steady-state probabilities indicate that 62.5%
of weeks will be in a bull market, 31.25% of weeks will be in a bear market and 6.25% of weeks
will be in a recession.
A thorough development and many examples can be found in the on-line monograph Meyn &
Tweedie 2005.The appendix of Meyn 2007, also available on-line, contains an abridged Meyn
& Tweedie.
Most of decision theory is normative or prescriptive, i.e., it is concerned with identifying the
best decision to take (in practice, there are situations in which "best" is not necessarily the
maximal (optimum may also include values in addition to maximum), but within a specific or
approximative range), assuming an ideal decision maker who is fully informed, able to compute
with perfect accuracy, and fully rational. The practical application of this prescriptive approach
(how people ought to make decisions) is called decision analysis, and aimed at finding tools,
methodologies and software to help people make better decisions. The most systematic and
comprehensive software tools developed in this way are called decision support systems.
Since people usually do not behave in ways consistent with axiomatic rules, often their own,
leading to violations of optimality, there is a related area of study, called a positive or
descriptive discipline, attempting to describe what people will actually do. Since the normative,
252
optimal decision often creates hypotheses for testing against actual behaviour, the two fields are
closely linked. Furthermore it is possible to relax the assumptions of perfect information,
rationality and so forth in various ways, and produce a series of different prescriptions or
predictions about behaviour, allowing for further tests of the kind of decision-making that
occurs in practice.
In recent decades, there has been increasing interest in what is sometimes called 'behavioral
decision theory' and this has contributed to a re-evaluation of what rational decision-making
requires (for instance Anand, 1993).
The revival of subjective probability theory, from the work of Frank Ramsey, Bruno de Finetti,
Leonard Savage and others, extended the scope of expected utility theory to situations where
subjective probabilities can be used. At this time, von Neumann's theory of expected utility
proved that expected utility maximization followed from basic postulates about rational
behavior.
The work of Maurice Allais and Daniel Ellsberg showed that human behavior has systematic
and sometimes important departures from expected-utility maximization. The prospect theory of
Daniel Kahneman and Amos Tversky renewed the empirical study of economic behavior with
less emphasis on rationality presuppositions. Kahneman and Tversky found three regularities —
in actual human decision-making, "losses loom larger than gains"; persons focus more on
changes in their utility–states than they focus on absolute utilities; and the estimation of
subjective probabilities is severely biased by anchoring.
General criticism
-Obs. The discussion above induces to consider the following class of systems: A system is
said to be PROBABILISTICALLY PREVISIBLE if it's possible to know all its outcomes, as well as
the correspondent probabilities;
253
Alternatives to decision theory
A highly controversial issue is whether one can replace the use of probability in decision theory
by other alternatives.
Probability theory
the work of Richard Threlkeld Cox for justification of the probability axioms,
the complete class theorems, which show that all admissible decision rules are
equivalent to the Bayesian decision rule for some utility function and some prior
distribution (or for the limit of a sequence of prior distributions). Thus, for every
decision rule, either the rule may be reformulated as a Bayesian procedure, or there is
a (perhaps limiting) Bayesian rule that is sometimes better and never worse.
The proponents of fuzzy logic, possibility theory, Dempster–Shafer theory and info-gap
decision theory maintain that probability is only one of many alternatives and point to many
examples where non-standard alternatives have been implemented with apparent success;
notably, probabilistic decision theory is sensitive to assumptions about the probabilities of
various events, while non-probabilistic rules such as minimax are robust, in that they do not
make such assumptions.
254
80- NO SÉCULO III AC, O MATEMÁTICO ITALIANO PINGALA, INVENTOU O “SISTEMA
BINÁRIO” DE NUMERAÇÃO; NELE, TODO NÚMERO NATURAL PODE SER ESCRITO COMO SOMA
DE POTÊNCIAS DE DOIS, MULTIPLICADAS PELOS COEFICIENTES “ZERO OU UM” (POR
EXEMPLO, 1 × 2³ + 0 × 2² + 1 × 21 + 1 × 20 = 11, OU SEJA, PODEMOS REPRESENTAR O
NÚMERO 11 PELA SEQUÊNCIA 1011 COMPOSTA PELOS NÚMEROS ZERO E UM); PARA
REPRESENTAR UM NÚMERO SUFICIENTE DE NÚMEROS, OPERAÇÕES E FUNÇÕES NUMÉRICAS,
LETRAS, PALAVRAS, SÍMBOLOS, ETC., USAM-SE SEQUÊNCIAS COMPOSTAS PELOS NÚMEROS
ZERO E UM (QUE COMPÕEM A CÓDIGO DE MÁQUINA).
EM COMPUTAÇÃO, UM DÍGITO BINÁRIO ( 0 OU 1 ) É CHAMADO DE bit (binary digit), E UM
AGRUPAMENTO DE 8 bits, CORRESPONDE A UM byte (binary term);
- SÉC. XVII, J.NAPIER (CUJO PRINCÍPIO DE ATRAVÉS DOS LOGARITMOS, TRANSFORMAR
PRODUTOS EM SOMAS, FOI APROVEITADO NAS RÉGUAS DE CÁLCULO), W. SCHICKARD
(EDIIN ?) (1623), B.PASCAL (1642) e LEIBNIZ (1670), PROJETARAM MÁQUINAS DE CALCULAR;
- EM 1703, G.LEIBNIZ DESENVOLVEU A LÓGICA, NUMA LINHA FORMAL E MATEMÁTICA,
USANDO O SISTEMA BINÁRIO (POR EXEMPLO: VERDADEIRO, 1 , e FALSO, 0 );
- EM 1765, O ENGENHEIRO E MATEMÁTICO ESCOCÊS JAMES WATT, INVENTOU UMA
MÁQUINA A VAPOR, CONSIDERADA UM MARCO PARA A REVOLUÇÃO INDUSTRIAL;
- EM 1801, O MECÂNICO FRANCÊS J.M.JAQUARD, INVENTOU UM TEAR MECÂNICO COM UMA
LEITORA AUTOMÁTICA DE CARTÕES PARA PRODUZIR DIFERENTES PADRÕES DE CORES;
- EM 1837, O ENGENHEIRO MECÂNICO E CIENTISTA INGLÊS, C. BABBAGE IDEALIZOU O
COMPUTADOR (CALCULADOR ANALÍTICO) COMO UM DISPOSITIVO PROGRAMÁVEL (ATRAVÉS
DE CARTÕES PERFURADOS) DE CÁLCULO, QUE NÃO CHEGOU A CONSTRUIR; COMO UM
COMPUTADOR MODERNO, TERIA UM PROCESSADOR PARA OS CÁLCULOS ARITMÉTICOS NO
SISTEMA DECIMAL, MEMÓRIA PARA REGISTRAR OS NÚMEROS, E A CAPACIDADE DE ALTERAR
SUA FUNÇÃO ATRAVÉS DE COMANDOS DO USUÁRIO, NO CASO, CARTÕES PERFURADOS;
- ADA LOVELACE, A PRIMEIRA PROGRAMADORA, INVENTOU A PRIMEIRA LINGUAGEM DE
COMPUTADOR E PUBLICOU OS PRIMEIROS PROGRAMAS, A PEDIDO DE BABBAGE;
- EM 1854, G.BOOLE PUBLICOU A SUA ÁLGEBRA BOOLEANA (REGRA ALGÉBRICAS QUE
TRADUZEM O ESSENCIAL DAS OPERAÇÕES LÓGICAS E , OU E NÃO , E DAS OPERAÇÕES DE
SOMA, PRODUTO E COMPLEMENTO DE CONJUNTOS;
- NO FINAL DO SÉC. XVIII, O NORTE-AMERICANO H.HOLLERITH, INVENTOU UMA MÁQUINA
CAPAZ DE PROCESSAR DADOS BASEADA NA SEPARAÇÃO ELETRÔNICA, PELOS FUROS DE
CARTÕES MAGNÉTICOS PERFURADOS;
- EM 1936/41, O ENGENHEIRO ALEMÃO K.ZUSE, CONSTRUIU O PRIMEIRO COMPUTADOR
PROGRAMÁVEL (ATRAVÉS DE FITAS PERFURADAS), COM MEMÓRIA, ELETRO-MECÂNICO
(UTILIZANDO RELÊS), E O SISTEMA BINÁRIO;
- EM 1945, K.ZUSE PROJETOU A 1ª LINGUAGEM DE ALTO NÍVEL, O PLANKALKÜL;
- EM 1936, A.TURING INVENTOU A “MÁQUINA DE TURING”: UMA MÁQUINA IDEAL, CAPAZ
DE REALIZAR TODO CÁLCULO (OU PROGRAMA) QUE UM COMPUTADOR DIGITAL PODE
REALIZAR;
- EM 1937, C.E.SHANNON APLICOU A ÁLGEBRA BOOLEANA EM CIRCUITOS
ELETRO-MECÂNICOS, PARA RESOLVER PROBLEMAS;
- DURANTE A 2ª GRANDE GUERRA, EM PROJETOS MILITARES:
- DE 1939/44, H. AIKEN, PROFESSOR DE HARVARD, PROJETOU O HARVARD MARK I,
COMPUTADOR ELETRO-MECÂNICO DE PROPÓSITO GERAL, BASEADO NO CALCULADOR
ANALÍTICO DE BABBAGE, COM VÁLVULAS NO LUGAR DOS DISPOSITIVOS MECÂNICOS DAS
CALCULADORAS, UTILIZANDO O SISTEMA DECIMAL, E CONSEGUINDO MULTIPLICAR DOIS
NÚMEROS DE 10 DÍGITOS EM 3 SEGUNDOS;
- DE 1941/46, OS ENGENHEIROS J.P.ECKERT e J. MAUCHY COORDENARAM O PROJETO DO
ENIAC, COMPUTADOR ELETRÔNICO DE PROPÓSITO GERAL, COM VÁLVULAS, UTILIZANDO O
SISTEMA DECIMAL, CAPAZ DE REALIZAR 100.000 CÁLCULOS SIMPLES OU QUINHENTAS
MULTIPLICAÇÕES, POR SEGUNDO, MAS QUE SÓ DEPOIS DE ALGUNS ANOS, SE TORNOU
255
PRATICAMENTE PROGRAMÁVEL; DE MODO DIFERENTE DO ENIAC, QUE UTILIZAVA
PROCESSAMENTO PARALELO, SEU SUCESSOR EDVAC, POSSUIA UMA ÚNICA UNIDADE DE
PROCESSAMENTO;
-JOHN VON NEWMANN, MATEMÁTICO HÚNGARO/ESTADUNIDENSE, SUGERIU O QUE VEIO A
SER CHAMADA DE ARQUITETURA DE VON NEWMANN,
1.Codificar as instruções de uma forma possível de ser armazenada na memória do
computador. Von Neumann sugeriu que fossem usados uns e zeros,
2.Armazenar as instruções na memória, bem como toda e qualquer informação necessária a
execução da tarefa, e
3.Quando processar o programa, buscar as instruções diretamente na memória, ao invés de
lerem um novo cartão perfurado a cada passo.
UTILIZADA NOS PROJETOS:
-EDVAC, DESENVOLVIDO PELA BALLISTIC RESEARCH LABORATORY U.S.,UNIVERSIDADE DA
PENSILVÂNIA, UTILIZANDO O SISTEMA BINÁRIO, E CONSIDERADO UM PROJETO DE SUCESSO;
MANCHESTER MARK I/ EDSAC (INSPIRADO NOS PLANOS DO EDVAC), BRITÂNICOS DE 1949:
256
UM ÚNICO CHIP; .
257
Frm Wikipedia. An INDUSTRIAL ROBOT is defined by ISO as an automatically
controlled, reprogrammable, multipurpose manipulator programmable in three or more
axes. The field of robotics may be more practically defined as the study, design and use
of robot systems for manufacturing (a top-level definition relying on the prior definition
of robot). Typical applications of robots include welding, painting, assembly, pick and
place (such as packaging, palletizing and SMT), product inspection, and testing; all
accomplished with high endurance, speed, and precision.
The most commonly used robot configurations are articulated robots, SCARA robots,
Delta robots and Cartesian coordinate robots, (aka gantry robots or x-y-z robots). In the
context of general robotics, most types of robots would fall into the category of robotic
arms (inherent in the use of the word manipulator in the above-mentioned ISO
standard). Robots exhibit varying degrees of autonomy:
Some robots are programmed to faithfully carry out specific actions over and over
again (repetitive actions) without variation and with a high degree of accuracy. These
actions are determined by programmed routines that specify the direction,
acceleration, velocity, deceleration, and distance of a series of coordinated motions.
Other robots are much more flexible as to the orientation of the object on which they
are operating or even the task that has to be performed on the object itself, which the
robot may even need to identify. For example, for more precise guidance, robots often
contain machine vision sub-systems acting as their "eyes", linked to powerful
computers or controllers. Artificial intelligence, or what passes for it, is becoming an
increasingly important factor in the modern industrial robot.
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A set of six-axis robots used for welding.
Factory Automation with industrial robots for palletizing food products like bread and
toast at a bakery in Germany
George Devol applied for the first robotics patents in 1954 (granted in 1961). The first company to
produce a robot was Unimation, founded by Devol and Joseph F. Engelberger in 1956, and was based on
Devol's original patents. Unimation robots were also called programmable transfer machines since their
main use at first was to transfer objects from one point to another, less than a dozen feet or so apart. They
used hydraulic actuators and were programmed in joint coordinates, i.e. the angles of the various joints
were stored during a teaching phase and replayed in operation. They were accurate to within 1/10,000 of
an inch[citation needed] (note: although accuracy is not an appropriate measure for robots, usually evaluated in
terms of repeatability - see later). Unimation later licensed their technology to Kawasaki Heavy Industries
and GKN, manufacturing Unimates in Japan and England respectively. For some time Unimation's only
competitor was Cincinnati Milacron Inc. of Ohio. This changed radically in the late 1970s when several
big Japanese conglomerates began producing similar industrial robots.
In 1969 Victor Scheinman at Stanford University invented the Stanford arm, an all-electric, 6-axis
articulated robot designed to permit an arm solution. This allowed it accurately to follow arbitrary paths
in space and widened the potential use of the robot to more sophisticated applications such as assembly
and welding. Scheinman then designed a second arm for the MIT AI Lab, called the "MIT arm."
Scheinman, after receiving a fellowship from Unimation to develop his designs, sold those designs to
Unimation who further developed them with support from General Motors and later marketed it as the
Programmable Universal Machine for Assembly (PUMA).
Industrial robotics took off quite quickly in Europe, with both ABB Robotics and KUKA Robotics
bringing robots to the market in 1973. ABB Robotics (formerly ASEA) introduced IRB 6, among the
world's first commercially available all electric micro-processor controlled robot. The first two IRB 6
robots were sold to Magnusson in Sweden for grinding and polishing pipe bends and were installed in
production in January 1974. Also in 1973 KUKA Robotics built its first robot, known as FAMULUS,[2]
also one of the first articulated robot to have six electromechanically driven axes.
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Interest in robotics increased in the late 1970s and many US companies entered the field, including large
firms like General Electric, and General Motors (which formed joint venture FANUC Robotics with
FANUC LTD of Japan). U.S. startup companies included Automatix and Adept Technology, Inc. At the
height of the robot boom in 1984, Unimation was acquired by Westinghouse Electric Corporation for 107
million U.S. dollars. Westinghouse sold Unimation to Stäubli Faverges SCA of France in 1988, which is
still making articulated robots for general industrial and cleanroom applications and even bought the
robotic division of Bosch in late 2004.
Only a few non-Japanese companies ultimately managed to survive in this market, the major ones being
Adept Technology, Stäubli-Unimation, the Swedish-Swiss company ABB Asea Brown Boveri and the
German company KUKA Robotics.
As of 2005, the robotic arm business is approaching a mature state, where they can provide enough speed,
accuracy and ease of use for most of the applications. Vision guidance (aka machine vision) is bringing a
lot of flexibility to robotic cells. However, the end effector attached to a robot is often a simple
pneumatic, 2-position chuck. This does not allow the robotic cell to easily handle different parts, in
different orientations.
Hand-in-hand with increasing off-line programmed applications, robot calibration is becoming more and
more important in order to guarantee a good positioning accuracy.
Other developments include downsizing industrial arms for light industrial use such as production of
small products, sealing and dispensing, quality control, handling samples in the laboratory. Such robots
are usually classified as "bench top" robots. Robots are used in pharmaceutical research in a technique
called High-throughput screening. Bench top robots are also used in consumer applications (micro-robotic
arms). Industrial arms may be used in combination with or even mounted on automated guided vehicles
(AGVs) to make the automation chain more flexible between pick-up and drop-off.
Market structure
According to the World Industrial Robotics 2011 report, there were 1,035,000 operational industrial
robots by the end of 2010. This number is estimated to reach 1,308,000 by the end of 2014.
The Japanese government estimates the industry could surge from about $5.2 billion in 2006 to $26
billion in 2010 and nearly $70 billion by 2025. In 2005, there were over 370,000 operational industrial
robots in Japan. A 2007 national technology roadmap by the Trade Ministry calls for 1 million industrial
robots to be installed throughout the country by 2025.
In August 2011, China Business News quoted Foxconn Chairman Terry Gou as saying the company
planned to use 1 million robots within three years, up from about 10,000 robots in use now and an
expected 300,000 next year.
Pricing Ranges: Typical stand-alone robot arms with welding packages cost between $28,000 and
$40,000. A pre-engineered workcell with safety equipment starts at $50,000 (in May 2012);
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The Reconditioned Option: Our top-quality refurbished robots offer the most cost-effective option. Used
robots cost 50% less than new models. Plus, you can rely on our exclusive reconditioning process to
produce reliable products.
Year Supply
1998 69,000
1999 79,000
2000 99,000
2001 78,000
2002 69,000
2003 81,000
2004 97,000
2005 120,000
2006 112,000
2007 114,000
2008 113,000
2009 60,000
2010 118,000
2011 150,000
De acordo com dados estatísticos publicados, o Japão possui 306 robôs/ 10 mil trabalhadores e
o Brasil 10 robôs/10 mil trabalhadores, contra uma média de 51 robôs/ 10 mil trabalhadores.
Um robô faz o serviço repetitivo de, no mínimo 3 trabalhadores/turno = 9 trabalhadores/dia.
Sought: an element x0 in A such that f(x0) ≤ f(x) for all x in A ("minimization") or such that
f(x0) ≥ f(x) for all x in A ("maximization").
Many real-world and theoretical problems may be modeled in this general framework.
Problems formulated using this technique in the fields of physics and computer vision may
refer to the technique as ENERGY MINIMIZATION, speaking of the value of the function f as
representing the ENERGY OF THE SYSTEM BEING MODELED.
The function f is called, variously, an objective function, cost function (minimization), utility
function (maximization), or, in certain fields, energy function, or energy functional. A feasible
solution that minimizes (or maximizes, if that is the goal) the objective function is called an
optimal solution.
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Fermat and Lagrange found calculus-based formulas for identifying optima, while Newton and
Gauss proposed iterative methods for moving towards an optimum;
Graph of a paraboloid given by f(x,y) = -(x²+y²)+4. The global maximum at (0,0,4) is indicated by
a red dot.
the method of LAGRANGE MULTIPLIERS (named after Joseph Louis Lagrange) provides a
strategy for finding the local maxima and minima of a function subject to equality constraints.
It is a powerful tool for solving this class of problems without the need to explicitly solve the
conditions and use them to eliminate extra variables.
Maximize
subject to
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Suppose we walk along the contour line with . In general the contour lines of
and may be distinct, so following the contour line for one could intersect
with or cross the contour lines of . This is equivalent to saying that while moving
along the contour line for the value of can vary. Only when the contour line
for meets contour lines of tangentially, do we not increase or decrease the
value of — that is, when the contour lines touch but do not cross.
The contour lines of f and g touch when the tangent vectors of the contour lines are
parallel. Since the gradient of a function is perpendicular to the contour lines, this is the
same as saying that the gradients of f and g are parallel. Thus we want points
where and
where
and
are the respective gradients. The constant is required because although the two
gradient vectors are parallel, the magnitudes of the gradient vectors are generally not
equal.
and solve
The constrained extrema of are critical points of the Lagrangian , but they are not
local extrema of (see Example 2 below).
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One may reformulate the Lagrangian as a Hamiltonian, in which case the solutions are
local minima for the Hamiltonian. This is done in optimal control theory, in the form of
Pontryagin's minimum principle.
The fact that solutions of the Lagrangian are not necessarily extrema also poses
difficulties for numerical optimization. This can be addressed by computing the
magnitude of the gradient, as the zeros of the magnitude are necessarily local minima,
as illustrated in the numerical optimization example.
Convex functions play an important role in many areas of mathematics. They are
especially important in the study of optimization problems where they are distinguished
by a number of convenient properties. For instance, a (strictly) convex function on an
open set has no more than one minimum. Even in infinite-dimensional spaces, under
suitable additional hypotheses, convex functions continue to satisfy such properties and,
as a result, they are the most well-understood functionals in the calculus of variations.
In probability theory, a convex function applied to the expected value of a random
variable is always less or equal to the expected value of the convex function of the
random variable. This result, known as Jensen's inequality underlies many important
inequalities (including, for instance, the arithmetic-geometric mean inequality and
Hölder's inequality).
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Formal definition: A real valued function f : X → R defined on a convex set X in a vector space
is called convex if, for any two points and in X and any ,
(note that is the slope of the purple line in the above drawing; note also
for all and in C. This condition is only slightly weaker than convexity. For
example, a real valued Lebesgue measurable function that is midpoint convex will be
convex. In particular, a continuous function that is midpoint convex will be convex.
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for all x and y in the interval. In particular, if f '(c) = 0, then c is a global minimum of
f(x).
A twice differentiable function of one variable is convex on an interval if and only if its
second derivative is non-negative there; this gives a practical test for convexity. If its
second derivative is positive then it is strictly convex, but the converse does not hold.
For example, the second derivative of f(x) = x4 is f "(x) = 12 x2, which is zero for x = 0,
but x4 is strictly convex.
Any local minimum of a convex function is also a global minimum. A strictly convex
function will have at most one global minimum.
For a convex function f, the sublevel sets {x | f(x) < a} and {x | f(x) ≤ a} with a ∈ R are
convex sets. However, a function whose sublevel sets are convex sets may fail to be a
convex function. A function whose sublevel sets are convex is called a quasiconvex
function.
for every
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defined on a convex subset of , the problem is to find any point in for
which the number is smallest, i.e., a point such that
for all .
The convexity of makes the powerful tools of convex analysis applicable. In finite-
dimensional normed spaces, the Hahn–Banach theorem and the existence of
subgradients lead to a particularly satisfying theory of necessary and sufficient
conditions for optimality, a duality theory generalizing that for linear programming, and
effective computational methods.
Standard form is the usual and most intuitive form of describing a convex minimization
problem. It consists of the following three parts:
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purposes, equality constraints are redundant; however, it can be beneficial to treat them
specially in practice.
The following statements are true about the convex minimization problem:
These results are used by the theory of convex minimization along with geometric
notions from functional analysis (in Hilbert spaces) such as the Hilbert projection
theorem, the separating hyperplane theorem, and Farkas' lemma.
Lagrange multipliers
For each point x in X that minimizes f over X, there exist real numbers λ0, ..., λm, called
Lagrange multipliers, that satisfy these conditions simultaneously:
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Conversely, if some x in X satisfies 1-3 for scalars λ0, ..., λm with λ0 = 1, then x is certain
to minimize f over X.
Methods
Cutting-plane methods
Ellipsoid method
Subgradient method
where x represents the vector of variables (to be determined), c and b are vectors of
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vectors are comparable when they have the same dimensions. If every entry in the first
is less-than or equal-to the corresponding entry in the second then we can say the first
vector is less-than or equal-to the second vector.
The problem of solving a system of linear inequalities dates back at least as far as
Fourier, after whom the method of Fourier-Motzkin elimination is named. The earliest
linear programming was first developed by Leonid Kantorovich in 1939: Kantorovich
developed the earliest linear programming problems in 1939 for use during World War
II to plan expenditures and returns in order to reduce costs to the army and increase
losses to the enemy. The method was kept secret until 1947 when George B. Dantzig
published the simplex method and John von Neumann developed the theory of duality
as a linear optimization solution, and applied it in the field of game theory. Postwar,
many industries found its use in their daily planning.
Dantzig's original example was to find the best assignment of 70 people to 70 jobs. The
computing power required to test all the permutations to select the best assignment is
vast; the number of possible configurations exceeds the number of particles in the
universe. However, it takes only a moment to find the optimum solution by posing the
problem as a linear program and applying the Simplex algorithm. The theory behind
linear programming drastically reduces the number of possible optimal solutions that
must be checked.
Linear programming can be applied to various fields of study. It is used in business and
economics, but can also be utilized for some engineering problems. Industries that use
linear programming models include transportation, energy, telecommunications, and
manufacturing. It has proved useful in modeling diverse types of problems in planning,
routing, scheduling, assignment, and design.
Interior point methods (also referred to as barrier methods) are a certain class of
algorithms to solve linear and nonlinear convex optimization problems.
Example solution
The interior point method was invented by John von Neumann. Von Neumann
suggested a new method of linear programming, using the homogeneous linear system
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of Gordan (1873) which was later popularized by Karmarkar's algorithm, developed by
Narendra Karmarkar in 1984 for linear programming. The method consists of a self-
concordant barrier function used to encode the convex set. Contrary to the simplex
method, it reaches an optimal solution by traversing the interior of the feasible region.
Any convex optimization problem can be transformed into minimizing (or maximizing)
a linear function over a convex set. The idea of encoding the feasible set using a barrier
and designing barrier methods was studied in the early 1960s by, amongst others,
Anthony V. Fiacco and Garth P. McCormick. These ideas were mainly developed for
general nonlinear programming, but they were later abandoned due to the presence of
more competitive methods for this class of problems (e.g. sequential quadratic
programming).
Yurii Nesterov and Arkadii Nemirovskii came up with a special class of such barriers
that can be used to encode any convex set. They guarantee that the number of iterations
of the algorithm is bounded by a polynomial in the dimension and accuracy of the
solution.
Karmarkar's breakthrough revitalized the study of interior point methods and barrier
problems, showing that it was possible to create an algorithm for linear programming
characterized by polynomial complexity and, moreover, that was competitive with the
simplex method. Already Khachiyan's ellipsoid method was a polynomial time
algorithm; however, in practice it was too slow to be of practical interest.
The class of primal-dual path-following interior point methods is considered the most
successful. Mehrotra's predictor-corrector algorithm provides the basis for most
implementations of this class of methods.
The usual objective of control theory is to calculate solutions for the proper corrective action
from the controller that result in system stability, that is, the system will hold the set point and
not oscillate around it.
The inputs and outputs of a continuous control system are generally related by nonlinear
differential equations. A transfer function can sometimes be obtained by:
The transfer function is also known as the system function or network function. The transfer
function is a mathematical representation, in terms of spatial or temporal frequency, of the
relation between the input and output of a linear time-invariant solution of the nonlinear
differential equations describing the system.
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Extensive use is usually made of a diagrammatic style known as the block diagram.
The concept of the feedback loop to control the dynamic behavior of the system: this is negative
feedback, because the sensed value is subtracted from the desired value to create the error
signal, which is amplified by the controller.
Control systems can be thought of as having four functions; Measure, Compare, Compute, and
Correct. These four functions are completed by five elements; Detector, Transducer,
Transmitter, Controller, and Final Control Element. The measuring function is completed by the
detector, transducer and transmitter. In practical applications these three elements are typically
contained in one unit. A standard example is a Resistance thermometer. The compare and
compute functions are completed within the controller which may be completed electronically
through a Proportional Control, PI Controller, PID Controller, Bistable, Hysteretic control or
Programmable logic controller. The correct function is completed with a final control element.
The final control element changes an input or output in the control system which affect the
manipulated or controlled variable.
- An example: Consider a car's cruise control, which is a device designed to maintain vehicle
speed at a constant desired or reference speed provided by the driver. The controller is the
cruise control, the plant is the car, and the system is the car and the cruise control. The system
output is the car's speed, and the control itself is the engine's throttle position which determines
how much power the engine generates.
A primitive way to implement cruise control is simply to lock the throttle position when the
driver engages cruise control. However, if the cruise control is engaged on a stretch of flat road,
then the car will travel slower going uphill and faster when going downhill. This type of
controller is called an open-loop controller because no measurement of the system output (the
car's speed) is used to alter the control (the throttle position.) As a result, the controller can not
compensate for changes acting on the car, like a change in the slope of the road.
In a closed-loop control system, a sensor monitors the system output (the car's speed) and feeds
the data to a controller which adjusts the control (the throttle position) as necessary to maintain
the desired system output (match the car's speed to the reference speed.) Now when the car goes
uphill the decrease in speed is measured, and the throttle position changed to increase engine
power, speeding the vehicle. Feedback from measuring the car's speed has allowed the
controller to dynamically compensate for changes to the car's speed. It is from this feedback that
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the paradigm of the control loop arises: the control affects the system output, which in turn is
measured and looped back to alter the control.
History
Although control systems of various types date back to antiquity, a more formal analysis of the
field began with a dynamics analysis of the centrifugal governor, conducted by the physicist
James Clerk Maxwell in 1868 entitled On Governors. This described and analyzed the
phenomenon of "hunting", in which lags in the system can lead to overcompensation and
unstable behavior. This generated a flurry of interest in the topic, during which Maxwell's
classmate Edward John Routh generalized Maxwell's results for the general class of linear
systems. Independently, Adolf Hurwitz analyzed system stability using differential equations in
1877, resulting in what is now known as the Routh–Hurwitz theorem.
A notable application of dynamic control was in the area of manned flight. The Wright brothers
made their first successful test flights on December 17, 1903 and were distinguished by their
ability to control their flights for substantial periods (more so than the ability to produce lift
from an airfoil, which was known). Continuous, reliable control of the airplane was necessary
for flights lasting longer than a few seconds.
By World War II, control theory was an important part of fire-control systems, guidance
systems and electronics.
Sometimes mechanical methods are used to improve the stability of systems. For example, ship
stabilizers are fins mounted beneath the waterline and emerging laterally. In contemporary
vessels, they may be gyroscopically controlled active fins, which have the capacity to change
their angle of attack to counteract roll caused by wind or waves acting on the ship.
The Sidewinder missile uses small control surfaces placed at the rear of the missile with
spinning disks on their outer surface; these are known as rollerons. Airflow over the disk spins
them to a high speed. If the missile starts to roll, the gyroscopic force of the disk drives the
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control surface into the airflow, cancelling the motion. Thus the Sidewinder team replaced a
potentially complex control system with a simple mechanical solution.
The Space Race also depended on accurate spacecraft control. However, control theory also saw
an increasing use in fields such as economics.
Many active and historical figures made significant contribution to control theory, including, for
example:
Alexander Lyapunov (1857–1918) in the 1890s marks the beginning of stability theory.
Harold S. Black (1898–1983), invented the concept of negative feedback amplifiers in
1927. He managed to develop stable negative feedback amplifiers in the 1930s.
Harry Nyquist (1889–1976), developed the Nyquist stability criterion for feedback
systems in the 1930s.
Richard Bellman (1920–1984), developed dynamic programming since the 1940s.
Andrey Kolmogorov (1903–1987) co-developed the Wiener–Kolmogorov filter (1941).
Norbert Wiener (1894–1964) co-developed the Wiener–Kolmogorov filter and coined
the term cybernetics in the 1940s.
John R. Ragazzini (1912–1988) introduced digital control and the z-transform in the
1950s.
Lev Pontryagin (1908–1988) introduced the maximum principle and the bang-bang
principle.
To avoid the problems of the open-loop controller, control theory introduces feedback. A
closed-loop controller uses feedback to control states or outputs of a dynamical system. Its name
comes from the information path in the system: process inputs (e.g., voltage applied to an
electric motor) have an effect on the process outputs (e.g., speed or torque of the motor), which
is measured with sensors and processed by the controller; the result (the control signal) is used
as input to the process, closing the loop.
In some systems, closed-loop and open-loop control are used simultaneously. In such systems,
the open-loop control is termed feedforward and serves to further improve reference tracking
performance.
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The output of the system y(t) is fed back through a sensor measurement F to the reference value
r(t). The controller C then takes the error e (difference) between the reference and the output to
change the inputs u to the system under control P. This is shown in the figure. This kind of
controller is a closed-loop controller or feedback controller.
If we assume the controller C, the plant P, and the sensor F are linear and time-invariant (i.e.,
elements of their transfer function C(s), P(s), and F(s) do not depend on time), the systems
above can be analysed using the Laplace transform on the variables. This gives the following
relations:
PID controller
The PID controller is probably the most-used feedback control design. PID is an acronym for
Proportional-Integral-Derivative, referring to the three terms operating on the error signal to
produce a control signal. If u(t) is the control signal sent to the system, y(t) is the measured
output and r(t) is the desired output, and tracking error , a PID
controller has the general form
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The desired closed loop dynamics is obtained by adjusting the three parameters , and
, often iteratively by "tuning" and without specific knowledge of a plant model. Stability
can often be ensured using only the proportional term. The integral term permits the rejection of
a step disturbance (often a striking specification in process control). The derivative term is used
to provide damping or shaping of the response. PID controllers are the most well established
class of control systems: however, they cannot be used in several more complicated cases,
especially if MIMO systems are considered.
In contrast to the frequency domain analysis of the classical control theory, modern control
theory utilizes the time-domain state space representation, a mathematical model of a physical
system as a set of input, output and state variables related by first-order differential equations.
To abstract from the number of inputs, outputs and states, the variables are expressed as vectors
and the differential and algebraic equations are written in matrix form (the latter only being
possible when the dynamical system is linear). The state space representation (also known as the
"time-domain approach") provides a convenient and compact way to model and analyze systems
with multiple inputs and outputs. With inputs and outputs, we would otherwise have to write
down Laplace transforms to encode all the information about a system. Unlike the frequency
domain approach, the use of the state space representation is not limited to systems with linear
components and zero initial conditions. "State space" refers to the space whose axes are the state
variables. The state of the system can be represented as a vector within that space.
Stability
The stability of a general dynamical system with no input can be described with Lyapunov
stability criteria. A linear system that takes an input is called bounded-input bounded-output
(BIBO) stable if its output will stay bounded for any bounded input. Stability for nonlinear
systems that take an input is input-to-state stability (ISS), which combines Lyapunov stability
and a notion similar to BIBO stability. For simplicity, the following descriptions focus on
continuous-time and discrete-time linear systems.
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Mathematically, this means that for a causal linear system to be stable all of the poles of its
transfer function must have negative-real values, i.e. the real part of all the poles are less than
zero. Practically speaking, stability requires that the transfer function complex poles reside
in the open left half of the complex plane for continuous time, when the Laplace
transform is used to obtain the transfer function.
inside the unit circle for discrete time, when the Z-transform is used.
The difference between the two cases is simply due to the traditional method of plotting
continuous time versus discrete time transfer functions. The continuous Laplace transform is in
Cartesian coordinates where the axis is the real axis and the discrete Z-transform is in
circular coordinates where the axis is the real axis.
When the appropriate conditions above are satisfied a system is said to be asymptotically stable:
the variables of an asymptotically stable control system always decrease from their initial value
and do not show permanent oscillations. Permanent oscillations occur when a pole has a real
part exactly equal to zero (in the continuous time case) or a modulus equal to one (in the
discrete time case). If a simply stable system response neither decays nor grows over time, and
has no oscillations, it is marginally stable: in this case the system transfer function has non-
repeated poles at complex plane origin (i.e. their real and complex component is zero in the
continuous time case). Oscillations are present when poles with real part equal to zero have an
imaginary part not equal to zero.
which has a pole in (zero imaginary part). This system is BIBO (asymptotically)
stable since the pole is inside the unit circle.
which has a pole at and is not BIBO stable since the pole has a modulus strictly
greater than one.
Numerous tools exist for the analysis of the poles of a system. These include graphical systems
like the root locus, Bode plots or the Nyquist plots.
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Mechanical changes can make equipment (and control systems) more stable. Sailors add ballast
to improve the stability of ships. Cruise ships use antiroll fins that extend transversely from the
side of the ship for perhaps 30 feet (10 m) and are continuously rotated about their axes to
develop forces that oppose the roll.
Controllability and observability are main issues in the analysis of a system before deciding the
best control strategy to be applied, or whether it is even possible to control or stabilize the
system. Controllability is related to the possibility of forcing the system into a particular state by
using an appropriate control signal. If a state is not controllable, then no signal will ever be able
to control the state. If a state is not controllable, but its dynamics are stable, then the state is
termed Stabilizable. Observability instead is related to the possibility of "observing", through
output measurements, the state of a system. If a state is not observable, the controller will never
be able to determine the behaviour of an unobservable state and hence cannot use it to stabilize
the system. However, similar to the stabilizability condition above, if a state cannot be observed
it might still be detectable.
From a geometrical point of view, looking at the states of each variable of the system to be
controlled, every "bad" state of these variables must be controllable and observable to ensure a
good behaviour in the closed-loop system. That is, if one of the eigenvalues of the system is not
both controllable and observable, this part of the dynamics will remain untouched in the closed-
loop system. If such an eigenvalue is not stable, the dynamics of this eigenvalue will be present
in the closed-loop system which therefore will be unstable. Unobservable poles are not present
in the transfer function realization of a state-space representation, which is why sometimes the
latter is preferred in dynamical systems analysis.
Control specification
Several different control strategies have been devised in the past years. These vary from
extremely general ones (PID controller), to others devoted to very particular classes of systems
(especially robotics or aircraft cruise control).
A control problem can have several specifications. Stability, of course, is always present: the
controller must ensure that the closed-loop system is stable, regardless of the open-loop
stability. A poor choice of controller can even worsen the stability of the open-loop system,
which must normally be avoided. Sometimes it would be desired to obtain particular dynamics
in the closed loop: i.e. that the poles have , where is a fixed value
strictly greater than zero, instead of simply asking that .
Another typical specification is the rejection of a step disturbance; including an integrator in the
open-loop chain (i.e. directly before the system under control) easily achieves this. Other classes
of disturbances need different types of sub-systems to be included.
Other "classical" control theory specifications regard the time-response of the closed-loop
system: these include the rise time (the time needed by the control system to reach the desired
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value after a perturbation), peak overshoot (the highest value reached by the response before
reaching the desired value) and others (settling time, quarter-decay). Frequency domain
specifications are usually related to robustness (see after).
A control system must always have some robustness property. A robust controller is such that
its properties do not change much if applied to a system slightly different from the mathematical
one used for its synthesis. This specification is important: no real physical system truly behaves
like the series of differential equations used to represent it mathematically. Typically a simpler
mathematical model is chosen in order to simplify calculations, otherwise the true system
dynamics can be so complicated that a complete model is impossible.
System identification
The process of determining the equations that govern the model's dynamics is called system
identification. This can be done off-line: for example, executing a series of measures from
which to calculate an approximated mathematical model, typically its transfer function or
matrix. Such identification from the output, however, cannot take account of unobservable
dynamics. Sometimes the model is built directly starting from known physical equations: for
example, in the case of a mass-spring-damper system we know that
. Even assuming that a "complete" model is used in
designing the controller, all the parameters included in these equations (called "nominal
parameters") are never known with absolute precision; the control system will have to behave
correctly even when connected to physical system with true parameter values away from
nominal.
Some advanced control techniques include an "on-line" identification process (see later). The
parameters of the model are calculated ("identified") while the controller itself is running: in this
way, if a drastic variation of the parameters ensues (for example, if the robot's arm releases a
weight), the controller will adjust itself consequently in order to ensure the correct performance.
Analysis
Analysis of the robustness of a SISO (single input single output) control system can be
performed in the frequency domain, considering the system's transfer function and using
Nyquist and Bode diagrams. Topics include gain and phase margin and amplitude margin. For
MIMO (multi input multi output) and, in general, more complicated control systems one must
consider the theoretical results devised for each control technique (see next section): i.e., if
particular robustness qualities are needed, the engineer must shift his attention to a control
technique by including them in its properties.
Constraints
A particular robustness issue is the requirement for a control system to perform properly in the
presence of input and state constraints. In the physical world every signal is limited. It could
279
happen that a controller will send control signals that cannot be followed by the physical
system: for example, trying to rotate a valve at excessive speed. This can produce undesired
behavior of the closed-loop system, or even damage or break actuators or other subsystems.
Specific control techniques are available to solve the problem: model predictive control (see
later), and anti-wind up systems. The latter consists of an additional control block that ensures
that the control signal never exceeds a given threshold.
System classifications
For MIMO systems, pole placement can be performed mathematically using a state space
representation of the open-loop system and calculating a feedback matrix assigning poles in the
desired positions. In complicated systems this can require computer-assisted calculation
capabilities, and cannot always ensure robustness. Furthermore, all system states are not in
general measured and so observers must be included and incorporated in pole placement design.
Processes in industries like robotics and the aerospace industry typically have strong nonlinear
dynamics. In control theory it is sometimes possible to linearize such classes of systems and
apply linear techniques, but in many cases it can be necessary to devise from scratch theories
permitting control of nonlinear systems. These, e.g., feedback linearization, backstepping,
sliding mode control, trajectory linearization control normally take advantage of results based
on Lyapunov's theory. Differential geometry has been widely used as a tool for generalizing
well-known linear control concepts to the non-linear case, as well as showing the subtleties that
make it a more challenging problem.
Decentralized systems
When the system is controlled by multiple controllers, the problem is one of decentralized
control. Decentralization is helpful in many ways, for instance, it helps control systems operate
over a larger geographical area. The agents in decentralized control systems can interact using
communication channels and coordinate their actions.
Every control system must guarantee first the stability of the closed-loop behavior. For linear
systems, this can be obtained by directly placing the poles. Non-linear control systems use
specific theories (normally based on Aleksandr Lyapunov's Theory) to ensure stability without
regard to the inner dynamics of the system. The possibility to fulfill different specifications
varies from the model considered and the control strategy chosen. Here a summary list of the
main control techniques is shown:
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Adaptive control
Hierarchical control
A Hierarchical control system is a type of Control System in which a set of devices and
governing software is arranged in a hierarchical tree. When the links in the tree are
implemented by a computer network, then that hierarchical control system is also a
form of Networked control system.
Intelligent control
Optimal control
Optimal control is a particular control technique in which the control signal optimizes a
certain "cost index": for example, in the case of a satellite, the jet thrusts needed to
bring it to desired trajectory that consume the least amount of fuel. Two optimal
control design methods have been widely used in industrial applications, as it has been
shown they can guarantee closed-loop stability. These are Model Predictive Control
(MPC) and linear-quadratic-Gaussian control (LQG). The first can more explicitly take
into account constraints on the signals in the system, which is an important feature in
many industrial processes. However, the "optimal control" structure in MPC is only a
means to achieve such a result, as it does not optimize a true performance index of the
closed-loop control system. Together with PID controllers, MPC systems are the most
widely used control technique in process control.
Robust control
Robust control deals explicitly with uncertainty in its approach to controller design.
Controllers designed using robust control methods tend to be able to cope with small
differences between the true system and the nominal model used for design. The early
methods of Bode and others were fairly robust; the state-space methods invented in
the 1960s and 1970s were sometimes found to lack robustness. A modern example of
a robust control technique is H-infinity loop-shaping developed by Duncan McFarlane
and Keith Glover of Cambridge University, United Kingdom. Robust methods aim to
achieve robust performance and/or stability in the presence of small modeling errors.
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Stochastic control
Stochastic control deals with control design with uncertainty in the model. In typical
stochastic control problems, it is assumed that there exist random noise and
disturbances in the model and the controller, and the control design must take into
account these random deviations.
Mathematical control theory is the area of application-oriented mathematics that deals with
the basic principles underlying the analysis and design of control systems. To control an object
means to influence its behavior so as to achieve a desired goal. In order to implement this
influence, engineers build devices that incorporate various mathematical techniques. These
devices range from Watt’s steam engine governor, designed during the English Industrial
Revolution, to the sophisticated microprocessor controllers found in consumer items —such as
CD players and automobiles— or in industrial robots and airplane autopilots.
The study of these devices and their interaction with the object being controlled is the subject
of this book. While on the one hand one wants to understand the fundamental limitations
that mathematics imposes on what is achievable, irrespective of the precise technology
being used, it is also true that technology may well influence the type of question to be
asked and the choice of mathematical model. An example of this is the use of difference
rather than differential equations when one is interested in digital control. Roughly speaking,
there have been two main lines of work in control theory, which sometimes have seemed to
proceed in very different directions but which are in fact complementary. One of these is
based on the idea that a good model of the object to be controlled is available and that one
wants to somehow optimize its behavior. For instance, physical principles and engineering
specifications can be —and are— used in order to calculate that trajectory of a spacecraft
which minimizes total travel time or fuel consumption. The techniques here are closely
related to the classical calculus of variations and to other areas of optimization theory; the
end result is typically a preprogrammed flight plan. The other main line of work is that based
on the constraints imposed by uncertainty about the model or about the environment in
which the object operates. The central tool here is the use of feedback in order to correct for
deviations from the desired behavior. For instance, various feedback control systems are used
during actual space flight in order to compensate for errors from the precomputed trajectory.
Mathematically, stability theory, dynamical systems, and especially the theory of functions of a
complex variable, have had a strong influence on this approach. It is widely recognized today
that these two broad lines of work deal just with different aspects of the same problems,
and we do not make an artificial distinction between them in this book. Later on we shall give
an axiomatic definition of what we mean by a “system” or “machine.” Its role will be
somewhat analogous to that played in mathematics by the definition of “function” as a set
of ordered pairs: not itself the object of study, but a necessary foundation upon which the
entire theoretical development will rest.
282
Romans, one finds water levels in aqueducts being kept constant through the use of various
combinations of valves. Modern developments started during the seventeenth century. The
Dutch mathematician and astronomer Christiaan Huygens designed pendulum clocks and in
doing so analyzed the problem of speed control; in this work he competed with his
contemporary Robert Hooke. The needs of navigation had a strong influence on scientific and
technological research at that time, and accurate clocks —to allow determinations of solar
time— were in great demand. The attention turned to windmills during the eighteenth
century, and speed controls based on Hooke’s and Huygens’ ideas were built. A central idea
here is the use of flyballs: Two balls attached to an axis turn together with the windmill, in
such a manner that centrifugal force due to angular velocity causes them to rise; in turn this
upward movement is made to affect the positions of the mill’s sails. Thus, feedback was
implemented by the linkages from the flyballs to the sails. But it was the Industrial Revolution,
and the adaptation by James Watt in 1769 of flyball governors to steam engines, that made
control mechanisms very popular; the problem there was to regulate engine speed despite a
variable load. Steady-state error could appear, and various inventors introduced variations of
the integral feedback idea in order to deal with this problem.
The mathematician and astronomer George Airy was the first to attempt, around 1840, an
analysis of the governor and similar devices. By 1868, there were about 75,000 Watt
governors in use; that year, the Scottish physicist James,Clerk Maxwell published the first
complete mathematical treatment of their properties and was able to explain the sometimes
erratic behavior that had been observed as well as the effect of integral control. His work gave
rise to the first wave of theoretical research in control, and characterizations of stability were
independently obtained for linear systems by the mathematicians A. Hurwitz and E.J. Routh.
This theory was applied in a variety of different areas, such as
the study of ship steering systems.
During the 1930s, researchers at Bell Telephone Laboratories developed the theory of
feedback amplifiers, based on assuring stability and appropriate response for electrical circuits.
This work, by H. Nyquist, H.W. Bode, and others, constitutes even today the foundation of
much of frequency design. Analog computers appeared also around that time and formed the
basis for implementing controllers in the chemical and petroleum industries. During the
Second World War, these techniques were used in the design of anti-aircraft batteries and fire-
control systems, applications that continue to be developed today. The mathematician
Norbert Wiener, who developed a theory of estimation for random processes used in these
applications, coined at the end of the war the term “cybernetics” to refer to control theory and
related areas. These so-called classical approaches were for the most part limited by their
restriction to linear time-invariant systems with scalar inputs and outputs. Only during the
1950s did control theory begin to develop powerful general techniques that allowed treating
multivariable, time-varying systems, as well as many nonlinear problems. Contributions by
Richard Bellman (dynamic programming) and Rudolf Kalman (filtering, linear/quadratic optimal
control, and algebraic analysis) in the United States, and by L. Pontryagin (nonlinear optimal
control) in the Soviet Union, formed the basis for a very large research effort during the
1960s, which continues to this day. Present day theoretical research in control theory involves
a variety of areas of pure mathematics. Concepts and results from these areas find applications
in control theory; conversely, questions about control systems give rise to new open problems
in mathematics.
283
activity, when this uncertainty is expressed in probabilistic or statistical terms. Mathematically
different but closely related is the area of robust control, which deals with the design of
control laws that are guaranteed to perform even if the assumed model of the system to be
controlled is incorrect —with the allowed deviations quantified in appropriate norms— or
under the possibility of imperfect controller design. The area of adaptive control deals also
with the control of partially unknown systems, but differs from robust control in the
mathematics employed. Adaptive controllers typically make use of identification techniques,
which produce estimates of the system parameters for use by controllers ;
When using computers one should consider the effect of quantization errors on the
implementation of control laws, which arise due to limited precision when real-valued signals
are translated into fixed-point representations (A/D or analog to digital conversion); see e.g.
[304]. Other questions relate to the interface between higher-level controllers implemented in
software and lowerlevel servos of an analog and physical character; this gives rise to the area
of hybrid systems.
Bang–bang control
284
Bang–bang solutions in optimal control
In optimal control problems, it is sometimes the case that a control is restricted to be between a
lower and an upper bound. If the optimal control switches from one extreme to the other at
certain times (i.e., is never strictly in between the bounds) then that control is referred to as a
bang-bang solution. Bang–bang controls frequently arise in minimum-time problems. For
example, if it is desired to stop a car in the shortest possible time at a certain position
sufficiently far ahead of the car, the solution is to apply maximum acceleration until the unique
switching point, and then apply maximum braking to come to rest exactly at the desired
position. This solution, which can be "uncomfortable" for the passengers, is a bang–bang
solution: maximum engine throttle followed by maximum braking.
A familiar everyday example is bringing water to a boil in the shortest time, which is achieved
by applying full heat, then turning it off when the water reaches a boil.
Bang–bang solutions also arise when the Hamiltonian is linear in the control variable;
application of Pontryagin's minimum principle will then lead to pushing the control to its upper
or lower bound depending on the sign of the coefficient of u in the Hamiltonian.
In summary, bang–bang controls are actually optimal controls in some cases, although they are
also often implemented because of simplicity or convenience.
Pontryagin's maximum (or minimum) principle is used in optimal control theory to find the
best possible control for taking a dynamical system from one state to another, especially in the
presence of constraints for the state or input controls. It was formulated in 1956 by the Russian
mathematician Lev Semenovich Pontryagin and his students. It has as a special case the Euler–
Lagrange equation of the calculus of variations.
The principle states informally that the Hamiltonian must be minimized over , the set of all
permissible controls. If is the optimal control for the problem, then the principle states
that:
The result was first successfully applied into minimum time problems where the input control is
constrained, but it can also be useful in studying state-constrained problems.
Special conditions for the Hamiltonian can also be derived. When the final time is fixed and
the Hamiltonian does not depend explicitly on time , then:
285
More general conditions on the optimal control are given below.
When satisfied along a trajectory, Pontryagin's minimum principle is a necessary condition for
an optimum. The Hamilton–Jacobi–Bellman equation provides sufficient conditions for an
optimum, but this condition must be satisfied over the whole of the state space.
The principle was first known as Pontryagin's maximum principle and its proof is historically
based on maximizing the Hamiltonian. The initial application of this principle was to the
maximization of the terminal velocity of a rocket. However as it was subsequently mostly used
for minimization of a performance index it has here been referred to as the minimum principle.
Pontryagin's book solved the problem of minimizing a performance index
When solved locally, the HJB is a necessary condition, but when solved over the whole
of state space, the HJB equation is a necessary and sufficient condition for an optimum.
The solution is open loop, but it also permits the solution of the closed loop problem.
The HJB method can be generalized to stochastic systems as well.
Classical variational problems, for example the brachistochrone problem, can be solved
using this method.
The equation is a result of the theory of dynamic programming which was pioneered in
the 1950s by Richard Bellman and coworkers.[1] The corresponding discrete-time
equation is usually referred to as the Bellman equation. In continuous time, the result
can be seen as an extension of earlier work in classical physics on the Hamilton-Jacobi
equation by William Rowan Hamilton and Carl Gustav Jacob Jacobi.
Consider the following problem in deterministic optimal control over the time period
:
where C[] is the scalar cost rate function and D[] is a function that gives the economic
value or utility at the final state, x(t) is the system state vector, x(0) is assumed given,
and u(t) for 0 ≤ t ≤ T is the control vector that we are trying to find.
286
where F[] gives the vector determining physical evolution of the state vector over time.
For this simple system, the Hamilton Jacobi Bellman partial differential equation is
where the means the dot product of the vectors a and b and is the gradient
operator.
The unknown scalar in the above PDE is the Bellman 'value function', which
represents the cost incurred from starting in state at time and controlling the system
optimally from then until time .
where o(dt2) denotes the terms in the Taylor expansion of higher order than one. Then if
we cancel V(x(t), t) on both sides, divide by dt, and take the limit as dt approaches zero,
we obtain the HJB equation defined above.
The HJB equation is usually solved backwards in time, starting from and ending
at .
When solved over the whole of state space, the HJB equation is a necessary and
sufficient condition for an optimum. If we can solve for then we can find from it a
control that achieves the minimum cost.
In general case, the HJB equation does not have a classical (smooth) solution. Several
notions of generalized solutions have been developed to cover such situations, including
287
viscosity solution (Pierre-Louis Lions and Michael Crandall), minimax solution (Andrei
Izmailovich Subbotin), and others.
The idea of solving a control problem by applying Bellman's principle of optimality and
then working out backwards in time an optimizing strategy can be generalized to
stochastic control problems. Consider similar as above
first using Bellman and then expanding with Itô's rule, one finds the
deterministic HJB equation
where represents the stochastic differentiation operator, and subject to the terminal
condition
Note, that the randomness has disappeared. In this case a solution of the latter does
not necessarily solve the primal problem, it is a candidate only and a further verifying
argument is required. This technique is widely used in Financial Mathematics to
determine optimal investment strategies in the market (see for example Merton's
portfolio problem).
Consider the following problem in deterministic optimal control over the time period
:
where C[] is the scalar cost rate function and D[] is a function that gives the economic
value or utility at the final state, x(t) is the system state vector, x(0) is assumed given,
and u(t) for 0 ≤ t ≤ T is the control vector that we are trying to find.
288
where F[] gives the vector determining physical evolution of the state vector over time.
For this simple system, the Hamilton Jacobi Bellman partial differential equation is
Where the means the dot product of the vectors a and b and is the gradient
operator.
The unknown scalar in the above PDE is the Bellman 'value function', which
represents the cost incurred from starting in state at time and controlling the system
optimally from then until time .
where o(dt2) denotes the terms in the Taylor expansion of higher order than one. Then if
we cancel V(x(t), t) on both sides, divide by dt, and take the limit as dt approaches zero,
we obtain the HJB equation defined above.
The HJB equation is usually solved backwards in time, starting from and ending
at .
When solved over the whole of state space, the HJB equation is a necessary and
sufficient condition for an optimum. [2] If we can solve for then we can find from it a
control that achieves the minimum cost.
In general case, the HJB equation does not have a classical (smooth) solution. Several
notions of generalized solutions have been developed to cover such situations, including
289
viscosity solution (Pierre-Louis Lions and Michael Crandall), minimax solution (Andrei
Izmailovich Subbotin), and others.
The idea of solving a control problem by applying Bellman's principle of optimality and
then working out backwards in time an optimizing strategy can be generalized to
stochastic control problems. Consider similar as above
first using Bellman and then expanding with Itô's rule, one finds the
deterministic HJB equation
where represents the stochastic differentiation operator, and subject to the terminal
condition
Note, that the randomness has disappeared. In this case a solution of the latter does
not necessarily solve the primal problem, it is a candidate only and a further verifying
argument is required. This technique is widely used in Financial Mathematics to
determine optimal investment strategies in the market (see for example Merton's
portfolio problem).
Pontryagin's minimum principle, necessary but not sufficient condition for optimum, by
minimizing a Hamiltonian, but this has the advantage over HJB of only needing to be
satisfied over the single trajectory being considered.
PERSPECTIVES
From Google: Mathematical Control Theory/
Unsolved Problems in Mathematical Systems and Control Theory
Edited by Vincent D. Blondel and Alexandre Megretski
Copyright c 2004 by Princeton University Press Published, PRINCETON AND
OXFORD.
290
were published in this volume attracted considerable attention in the
research community. Early versions of some of the problems in this book
have been presented at the Open Problem sessions of the Oberwolfach
Tagung on Regelungstheorie, on February 27, 2002, and of the Conference
on Mathematical Theory of Networks and Systems (MTNS) in Notre Dame,
Indiana, on August 12, 2002.
The editors thank the organizers of these meetings for their willingness to
provide the problems this welcome exposure.”
291
the purposes of this meeting and the principles on which it has been organized. (For a
detailed listing of the 31 plenary speakers and their topics, see the “Meetings” section of the
present issue of the Notices or the Web site http://www.ams.org/meetings/; for a discursive
description of the meeting, see Allyn Jackson’s article “Stellar Lineup for UCLA Meeting” in
the February 2000 Notices.)
The program for “Mathematical Challenges” reflects the way in which mathematics is
reaching out to other disciplines, solving problems, and opening new pathways for research,
while at the same time drawing in ideas that bring new vitality and richness to the field. The
speakers have been asked to give a broad picture of the prospects and challenges in the
areas they cover and to do so in terms that are comprehensible to a general mathematical
audience. This will be an important event, and I urge and welcome you to attend.
—Felix E. Browder
292
Stellar Lineup for AMS Meeting in August 2000
The AMS meeting in August 2000 is shaping up to be a landmark event. Entitled
Mathematical Challenges of the 21st Century, the meeting will feature about
thirty speakers chosen from the world’s top mathematicians. These mathematical
leaders will share their ideas about the most important questions in
mathematics today and what the future might bring. The meeting will be held
August 7–12, 2000, on the campus of the University of California, Los Angeles.
“Mathematical Challenges” rivals the International Congresses of
Mathematicians (ICMs) in the number and prominence of the invited speakers.
The elite group of speakers chosen for “Mathematical Challenges” includes eight
Fields Medalists and a number of others who are generally considered to have
been top contenders for this distinction. Within the group one also finds a
Nevanlinna prizewinner, a King Faisal prizewinner, recipients of the Wolf
Prize, a Turing Award winner, and several AMS prizewinners. But more
impressive than the prizes is the way this group of speakers exemplifies the depth
and range of modern mathematics. Several will address aspects of number
theory, particularly the “Langlands Program”, a web of conjectures linking
disparate parts of the subject. Other speakers will discuss progress on another
great theme in number theory, the Riemann Hypothesis. The program’s emphasis
on geometry and topology reflects the importance of these areas in mathematics
today. Some of the most exciting developments in mathematics in recent years
have involved connections between geometry and topology on the one hand and
theoretical physics on the other, and several speakers will address this theme.
Another theme in the program is the impact of mathematics in biology, where
the problems require highly sophisticated mathematical models and present
formidable computational challenges. Indeed, the problem of creating more
powerful computational methods in all areas of science and technology is a
major motivation in mathematics today. The program for “Mathematical
Challenges” will demonstrate the impact of mathematics across many areas of
human endeavor, from science to commerce to communications to medicine. The
speakers for “Mathematical Challenges” are being urged to make their talks
accessible to a broad mathematical audience. In addition, because the meeting is
designed to provide a look toward the future of mathematics, the speakers are
being encouraged to speculate on their ideas about where the field is headed and
what kinds of problems, themes, and ideas will come to the fore in the next century.
“Mathematical Challenges of the 21st Century” promises to be a meeting of
historical significance. Allyn Jackson
293
OBSERVAÇÕES:
4) O EXEMPLO DA PG. 320 FOI OBTIDO ATRAVÉS DO SITE mathoverflow, que parece
responder a professores universitários. Parece que o Mathematics Stack Exchanges é
extenso a alunos.
L2) TEORIA ABSTRATA EXEMPLOS DE APLICAÇÃO, uma vez que esta última
pressupõe
294
OBSERVAÇÕES:
295